An Introduction to the Mathematical Theory of Inverse Problems [3 ed.] 9783030633424, 9783030633431

Table of contents :
Preface to the Third Edition
Preface to the Second Edition
Preface to the First Edition
Contents
1 Introduction and Basic Concepts
1.1 Examples of Inverse Problems
1.2 Ill-Posed Problems
1.3 The Worst-Case Error
1.4 Problems
2 Regularization Theory for Equations of the First Kind
2.1 A General Regularization Theory
2.2 Tikhonov Regularization
2.3 Landweber Iteration
2.4 A Numerical Example
2.5 The Discrepancy Principle of Morozov
2.6 Landweber's Iteration Method with Stopping Rule
2.7 The Conjugate Gradient Method
2.8 Problems
3 Regularization by Discretization
3.1 Projection Methods
3.2 Galerkin Methods
3.2.1 The Least Squares Method
3.2.2 The Dual Least Squares Method
3.2.3 The Bubnov–Galerkin Method for Coercive Operators
3.3 Application to Symm's Integral Equation of the First Kind
3.4 Collocation Methods
3.4.1 Minimum Norm Collocation
3.4.2 Collocation of Symm's Equation
3.5 Numerical Experiments for Symm's Equation
3.6 The Backus–Gilbert Method
3.7 Problems
4 Nonlinear Inverse Problems
4.1 Local Illposedness
4.2 The Nonlinear Tikhonov Regularization
4.2.1 Existence of Solutions and Stability
4.2.2 Source Conditions And Convergence Rates
4.2.3 A Parameter-Identification Problem
4.2.4 A Glimpse on Extensions to Banach Spaces
4.3 The Nonlinear Landweber Iteration
4.4 Problems
5 Inverse Eigenvalue Problems
5.1 Introduction
5.2 Construction of a Fundamental System
5.3 Asymptotics of the Eigenvalues and Eigenfunctions
5.4 Some Hyperbolic Problems
5.5 The Inverse Problem
5.6 A Parameter Identification Problem
5.7 Numerical Reconstruction Techniques
5.8 Problems
6 An Inverse Problem in Electrical Impedance Tomography
6.1 Introduction
6.2 The Direct Problem and the Neumann–Dirichlet Operator
6.3 The Inverse Problem
6.4 The Factorization Method
6.5 Problems
7 An Inverse Scattering Problem
7.1 Introduction
7.2 The Direct Scattering Problem
7.3 Properties of the Far Field Patterns
7.4 Uniqueness of the Inverse Problem
7.5 The Factorization Method
7.6 The Interior Transmission Eigenvalue Problem
7.6.1 The Radially Symmetric Case
7.6.2 Discreteness And Existence in the General Case
7.6.3 The Inverse Spectral Problem for the Radially Symmetric Case
7.7 Numerical Methods
7.7.1 A Simplified Newton Method
7.7.2 A Modified Gradient Method
7.7.3 The Dual Space Method
7.8 Problems
A Basic Facts from Functional Analysis
A.1 Normed Spaces and Hilbert Spaces
A.2 Orthonormal Systems
A.3 Linear Bounded and Compact Operators
A.4 Sobolev Spaces of Periodic Functions
A.5 Sobolev Spaces on the Unit Disc
A.6 Spectral Theory for Compact Operators in Hilbert Spaces
A.7 The Fréchet Derivative
A.8 Convex Analysis
A.9 Weak Topologies
A.10 Problems
B Proofs of the Results of Section 2.7摥映數爠eflinksspscgsps22.72

Citation preview

Applied Mathematical Sciences

Andreas Kirsch

An Introduction to the Mathematical Theory of Inverse Problems Third Edition

Applied Mathematical Sciences Volume 120

Series Editors Anthony Bloch, Department of Mathematics, University of Michigan, Ann Arbor, MI, USA C. L. Epstein, Department of Mathematics, University of Pennsylvania, Philadelphia, PA, USA Alain Goriely, Department of Mathematics, University of Oxford, Oxford, UK Leslie Greengard, New York University, New York, NY, USA Advisory Editors J. Bell, Center for Computational Sciences and Engineering, Lawrence Berkeley National Laboratory, Berkeley, CA, USA P. Constantin, Department of Mathematics, Princeton University, Princeton, NJ, USA R. Durrett, Department of Mathematics, Duke University, Durham, CA, USA R. Kohn, Courant Institute of Mathematical Sciences, New York University, New York, NY, USA R. Pego, Department of Mathematical Sciences, Carnegie Mellon University, Pittsburgh, PA, USA L. Ryzhik, Department of Mathematics, Stanford University, Stanford, CA, USA A. Singer, Department of Mathematics, Princeton University, Princeton, NJ, USA A. Stevens, Department of Applied Mathematics, University of Münster, Münster, Germany S. Wright, Computer Sciences Department, University of Wisconsin, Madison, WI, USA Founding Editors F. John, New York University, New York, NY, USA J. P. LaSalle, Brown University, Providence, RI, USA L. Sirovich, Brown University, Providence, RI, USA

The mathematization of all sciences, the fading of traditional scientific boundaries, the impact of computer technology, the growing importance of computer modeling and the necessity of scientific planning all create the need both in education and research for books that are introductory to and abreast of these developments.The purpose of this series is to provide such books, suitable for the user of mathematics, the mathematician interested in applications, and the student scientist. In particular, this series will provide an outlet for topics of immediate interest because of the novelty of its treatment of an application or of mathematics being applied or lying close to applications. These books should be accessible to readers versed in mathematics or science and engineering, and will feature a lively tutorial style, a focus on topics of current interest, and present clear exposition of broad appeal. A compliment to the Applied Mathematical Sciences series is the Texts in Applied Mathematics series, which publishes textbooks suitable for advanced undergraduate and beginning graduate courses.

More information about this series at http://www.springer.com/series/34

Andreas Kirsch

An Introduction to the Mathematical Theory of Inverse Problems Third Edition

123

Andreas Kirsch Department of Mathematics Karlsruhe Institute of Technology (KIT) Karlsruhe, Germany

ISSN 0066-5452 ISSN 2196-968X (electronic) Applied Mathematical Sciences ISBN 978-3-030-63342-4 ISBN 978-3-030-63343-1 (eBook) https://doi.org/10.1007/978-3-030-63343-1 Mathematics Subject Classification: 45Q05, 47A52, 47J06, 65J20, 65J22, 65L08, 65L09, 65N20, 65N21, 65R30, 65R32, 78A45, 78A46, 81U40, 86A22, 31A25, 34A55, 34B24, 35J05, 35R25, 25R30 1st edition: © Springer-Verlag New York, Inc. 1996 2nd edition: © Springer Science+Business Media, LLC 2011 3rd edition: © Springer Nature Switzerland AG 2021 This work is subject to copyright. All rights are reserved by the Publisher, whether the whole or part of the material is concerned, specifically the rights of translation, reprinting, reuse of illustrations, recitation, broadcasting, reproduction on microfilms or in any other physical way, and transmission or information storage and retrieval, electronic adaptation, computer software, or by similar or dissimilar methodology now known or hereafter developed. The use of general descriptive names, registered names, trademarks, service marks, etc. in this publication does not imply, even in the absence of a specific statement, that such names are exempt from the relevant protective laws and regulations and therefore free for general use. The publisher, the authors and the editors are safe to assume that the advice and information in this book are believed to be true and accurate at the date of publication. Neither the publisher nor the authors or the editors give a warranty, expressed or implied, with respect to the material contained herein or for any errors or omissions that may have been made. The publisher remains neutral with regard to jurisdictional claims in published maps and institutional affiliations. This Springer imprint is published by the registered company Springer Nature Switzerland AG The registered company address is: Gewerbestrasse 11, 6330 Cham, Switzerland

Preface to the Third Edition The field of inverse problems is growing rapidly, and during the 9 years since the appearance of the second edition of this book many new aspects and subfields have been developed. Since, obviously, not every subject can be treated in a single monograph, the author had to make a decision. As I pointed out already in the preface of the first edition, my intention was—and still is—to introduce the reader to some of the basic principles and developments of this field of inverse problems rather than going too deeply into special topics. As I continued to lecture on inverse problems at the University of Karlsruhe (now Karlsruhe Institute of Technology), new material has been added to the courses and thus also to this new edition because the idea of this book is still to serve as a type of textbook for a course on inverse problems. I have decided to extend this monograph in two directions. For some readers, it was perhaps a little unsatisfactory that only the abstract theory for linear problems was presented but the applications to inverse eigenvalue problems, electrical impedance tomography, and inverse scattering theory are of a nonlinear type. For that reason, and also because the abstract theories for Tikhonov’s method and Landweber’s iteration for nonlinear equations have come to a certain completion, I included a new chapter (Chapter 4) on these techniques for locally improperly posed nonlinear equations in Hilbert spaces with an outlook into some rather new developments for Banach spaces. The former Chapters 4, 5, and 6 are moved to 5, 6, and 7, respectively. The additional functional analytic tools needed in this new Chapter 4 result in two new sections of Appendix A on convex analysis and weak topologies. As a second new topic, a separate section (Section 7.6) on interior transmission eigenvalues is included. These eigenvalue problems arise naturally in the study of inverse scattering problems for inhomogeneous media and were introduced already in the former editions of this monograph. Besides their importance in scattering theory, the transmission eigenvalue problem is an interesting subject in itself, mainly because it fails to be self-adjoint. The investigation of the spectrum is a subject of the present

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research. Special issues of Inverse Problems [37] and recent monographs [34, 55] have addressed this topic alreadyf for the study of complex eigenvalues, one is until now restricted to radially symmetric refractive indices which reduces the partial differential equations to ordinary differential equations. Classical tools from complex analysis make it possible to prove the existence of complex eigenvalues (Subsection 7.6.1) and uniqueness for the corresponding inverse spectral problem (Subsection 7.6.3). I think that this analogue to the inverse Sturm–Liouville problem of Chapter 5 is a natural completion of the study of interior transmission eigenvalues. Finally, a rather large number of mistakes, ambiguities, and misleading formulations has been corrected in every chapter. As major mistakes, the proofs of Theorems 4.22 (a) and 6.30 (d) (referring to the numbering of the second edition) have been corrected. I want to thank all of my colleagues and the readers of the first two editions for the overwhelming positive responses and, last but not least, the publisher for its encouragement for writing this third edition. Karlsruhe, Germany December 2020

Andreas Kirsch

Preface to the Second Edition The first edition of the book appeared 14 years ago. The area of inverse problems is still a growing field of applied mathematics and an attempt at a second edition after such a long time was a difficult task for me. The number of publications on the subjects treated in this book has grown considerably and a new generation of mathematicians, physicists, and engineers has brought new concepts into the field. My philosophy, however, has never been to present a comprehensive book on inverse problems that covers all aspects. My purpose was (as I pointed out in the preface of the first edition), and still is, to present a book that can serve as a basis for an introductory (graduate) course in this field. The choice of material covered in this book reects my personal point of view: students should learn the basic facts for linear ill-posed problems including some of the present classical concepts of regularization and also some important examples of more modern nonlinear inverse problems. Although there has been considerable progress made on regularization concepts and convergence properties of iterative methods for abstract nonlinear inverse problems, I decided not to include these new developments in this monograph. One reason is that these theoretical results on nonlinear inverse problems are still not applicable to the inverse scattering problems that are my major field of interest. Instead, I refer the reader to the monographs [92, 149] where regularization methods for nonlinear problems are intensively treated. Also, in my opinion, every nonlinear inverse problem has its own characteristic features that should be used for a successful solution. With respect to the inverse scattering problem to determine the shape of the support of the contrast, a whole class of methods has been developed during the last decade, sometimes subsumed under the name Sampling Methods. Because they are very popular not only in the field of inverse scattering theory but also in the field of electrical impedance tomography (EIT) I decided to include the Factorization Method as one of the prominent members in this monograph.

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The Factorization Method is particularly simple for the problem of EIT and this field has attracted a lot of attention during the past decade, therefore a chapter on EIT has been added to this monograph as Chapter 5 and the chapter on inverse scattering theory now becomes Chapter 6. The main changes of this second edition compared to the first edition concern only Chapters 5 and 6 and Appendix A. As just mentioned, in Chapter 5 we introduce the reader to the inverse problem of electrical impedance tomography. This area has become increasingly important because of its applications in medicine and engineering sciences. Also, the methods of EIT serve as tools and guidelines for the investigation of other areas of tomography such that optical and photoacoustic tomography. The forward model of EIT is usually set up in the weak sense, that is, in appropriate Sobolev spaces. Although I expect that the reader is familiar with the basic facts on Sobolev spaces such as the trace theorem and Friedrich’s inequality, a tutorial section on Sobolev spaces on the unit disk is added in Appendix A, Section A.5. The approach using Fourier techniques is not very common but fits well with the presentation of Sobolev spaces of fractional order on the boundary of the unit disk in Section A.4 of Appendix A. In Chapter 5 on electrical impedance tomography the Neumann–Dirichlet operator is introduced and its most important properties such as monotonicity, continuity, and differentiability are shown. Uniqueness of the inverse problem is proven for the linearized problem only because it was this example for which Calderón presented his famous proof of uniqueness. (The fairly recent uniqueness proof by Astala and Päivärinta in [10] is far too complicated to be treated in this introductory work.) As mentioned above, the Factorization Method was developed during the last decade. It is a completely new and mathematically elegant approach to characterize the shape of the domain where the conductivity differs from the background by the Neumann– Dirichlet operator. The Factorization Method is an example of an approach that uses special features of the nonlinear inverse problem under consideration and has no analogy for traditional linear inverse problems. Major changes are also made in Chapter 6 on inverse scattering problems. A section on the Factorization Method has been added (Section 6.4) because inverse scattering problems are the type of problem for which it is perfectly applicable. The rigorous mathematical treatment of the Factorization Method makes it necessary to work with weak solutions of the scattering problem. Therefore, here we also have to use (local) Sobolev spaces rather than spaces of continuously differentiable functions. I took the opportunity to introduce the reader to a (in my opinion) very natural approach to prove existence of weak solutions by the Lippmann–Schwinger equation in L2(D) (where D contains the support of the contrast n − 1). The key is the fact that the volume potential with any L2-density solves the corresponding inhomogeneous Helmholtz equation in the weak sense (just as in the case of smooth densities) and can easily be proved by using the classical result and a density argument. The notion of weak solutions has the advantage of allowing

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arbitrary L∞-functions as indices of refraction but makes it necessary to modify almost all of the arguments in this chapter slightly. In Section 6.7 we dropped the motivating example for the uniqueness of the inverse scattering problem (Lemma 6.8 in the first edition) because it has already been presented for the uniqueness of the linearized inverse problem of impedance tomography. Finally, I want to thank all the readers of the first edition of the monograph for their extraordinarily positive response. I hope that with this second edition I added some course material suitable for being presented in a graduate course on inverse problems. In particular I have found that my students like the problem of impedance tomography and, in particular, the Factorization Method and I hope that this is true for others! Karlsruhe, Germany March 2011

Andreas Kirsch

Preface to the First Edition Following Keller [152] we call two problems inverse to each other if the formulation of each of them requires full or partial knowledge of the other. By this definition, it is obviously arbitrary which of the two problems we call the direct and which we call the inverse problem. But usually, one of the problems has been studied earlier and, perhaps, in more detail. This one is usually called the direct problem, whereas the other is the inverse problem. However, there is often another more important difference between these two problems. Hadamard (see [115]) introduced the concept of a wellposed prob lem, originating from the philosophy that the mathematical model of a physical problem has to have the properties of uniqueness, existence, and stability of the solution. If one of the properties fails to hold, he called the problem ill-posed. It turns out that many interesting and important inverse problems in science lead to ill-posed problems, whereas the corresponding direct problems are well-posed. Often, existence and uniqueness can be forced by enlarging or reducing the solution space (the space of “models”). For restoring stability, however, one has to change the topology of the spaces, which is in many cases impossible because of the presence of measurement errors. At first glance, it seems to be impossible to compute the solution of a problem numerically if the solution of the problem does not depend continuously on the data, that is, for the case of ill-posed problems. Under additional a priori information about the solution, such as smoothness and bounds on the derivatives, however, it is possible to restore stability and construct efficient numerical algorithms. We make no claim to cover all of the topics in the theory of inverse problems. Indeed, with the rapid growth of this field and its relationship to many fields of natural and technical sciences, such a task would certainly be impossible for a single author in a single volume. The aim of this book is twofold: first, we introduce the reader to the basic notions and difficulties encountered with ill-posed problems. We then study the basic properties of regularization methods for linear ill-posed problems. These methods can roughly be classified into two groups, namely, whether the regularization

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parameter is chosen a priori or a posteriori. We study some of the most important regularization schemes in detail. The second aim of this book is to give a first insight into two special nonlinear inverse problems that are of vital importance in many areas of the applied sciences. In both inverse spectral theory and inverse scattering theory, one tries to determine a coefficient in a differential equation from measurements of either the eigenvalues of the problem or the field “far away” from the scatterer. We hope that these two examples clearly show that a successful treatment of nonlinear inverse problems requires a solid knowledge of characteristic features of the corresponding direct problem. The combination of classical analysis and modern areas of applied and numerical analysis is, in the author’s opinion, one of the fascinating features of this relatively new area of applied mathematics. This book arose from a number of graduate courses, lectures, and survey talks during my time at the universities of Göttingen and Erlangen/ Nürnberg. It was my intention to present a fairly elementary and complete introduction to the field of inverse problems, accessible not only to mathematicians but also to physicists and engineers. I tried to include as many proofs as possible as long as they required knowledge only of classical differential and integral calculus. The notions of functional analysis make it possible to treat different kinds of inverse problems in a common language and extract its basic features. For the convenience of the reader, I have collected the basic definitions and theorems from linear and nonlinear functional analysis at the end of the book in an appendix. Results on nonlinear mappings, in particular for the Fréchet derivative, are only needed in Chapters 4 and 5. The book is organized as follows. In Chapter 1, we begin with a list of pairs of direct and inverse problems. Many of them are quite elementary and should be well known. We formulate them from the point of view of inverse theory to demonstrate that the study of particular inverse problems has a long history. Sections 1.3 and 1.4 introduce the notions of ill-posedness and the worstcase error. Although ill-posedness of a problem (roughly speaking) implies that the solution cannot be computed numerically — which is a very pessimistic point of view — the notion of the worst-case error leads to the possibility that stability can be recovered if additional information is available. We illustrate these notions with several elementary examples. In Chapter 2, we study the general regularization theory for linear ill-posed equations in Hilbert spaces. The general concept in Section 2.1 is followed by the most important special examples: Tikhonov regularization in Section 2.2, Landweber iteration in Section 2.3, and spectral cutoff in Section 2.4. These regularization methods are applied to a test example in Section 2.5. While in Sections 2.1–2.5 the regularization parameter has been chosen a priori, that is before starting the actual computation, Sections 2.6– 2.8 are devoted to regularization methods in which the regularization parameter is chosen implicitly by the stopping rule of the algorithm. In

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Sections 2.6 and 2.7, we study Morozov’s discrepancy principle and, again, Landweber’s iteration method. In contrast to these linear regularization schemes, we will investigate the conjugate gradient method in Section 2.8. This algorithm can be interpreted as a nonlinear regularization method and is much more difficult to analyze. Chapter 2 deals with ill-posed problems in infinite-dimensional spaces. However, in practical situations, these problems are first discretized. The discretization of linear ill-posed problems leads to badly conditioned finite linear systems. This subject is treated in Chapter 3. In Section 3.1, we recall basic facts about general projection methods. In Section 3.2, we study several Galerkin methods as special cases and apply the results to Symm’s integral equation in Section 3.3. This equation serves as a popular model equation in many papers on the numerical treatment of integral equations of the first kind with weakly singular kernels. We present a complete and elementary existence and uniqueness theory of this equation in Sobolev spaces and apply the results about Galerkin methods to this equation. In Section 3.4, we study collocation methods. Here, we restrict ourselves to two examples: the moment collocation and the collocation of Symm’s integral equation with trigonometric polynomials or piecewise constant functions as basis functions. In Section 3.5, we compare the different regularization techniques for a concrete numerical example of Symm’s integral equation. Chapter 3 is completed by an investigation of the Backus–Gilbert method. Although this method does not quite fit into the general regularization theory, it is nevertheless widely used in the applied sciences to solve moment problems. In Chapter 4, we study an inverse eigenvalue problem for a linear ordinary differential equation of second order. In Sections 4.2 and 4.3, we develop a careful analysis of the direct problem, which includes the asymptotic behaviour of the eigenvalues and eigenfunctions. Section 4.4 is devoted to the question of uniqueness of the inverse problem, that is, the problem of recovering the coefficient in the differential equation from the knowledge of one or two spectra. In Section 4.5, we show that this inverse problem is closely related to a parameter identification problem for parabolic equations. Section 4.6 describes some numerical reconstruction techniques for the inverse spectral problem. In Chapter 5, we introduce the reader to the field of inverse scattering theory. Inverse scattering problems occur in several areas of science and technology, such as medical imaging, nondestructive testing of material, and geological prospecting. In Section 5.2, we study the direct problem and prove uniqueness, existence, and continuous dependence on the data. In Section 5.3, we study the asymptotic form of the scattered field as r ! 1 and introduce the far eld pattern. The corresponding inverse scattering problem is to recover the index of refraction from a knowledge of the far field pattern. We give a complete proof of uniqueness of this inverse problem in Section 5.4. Finally, Section 5.5 is devoted to the study of some recent reconstruction techniques for the inverse scattering problem.

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Chapter 5 differs from previous ones in the unavoidable fact that we have to use some results from scattering theory without giving proofs. We only formulate these results, and for the proofs we refer to easily accessible standard literature. There exists a tremendous amount of literature on several aspects of inverse theory ranging from abstract regularization concepts to very concrete applications. Instead of trying to give a complete list of all relevant contributions, I mention only the monographs [17, 105, 110, 136, 168, 173, 174, 175, 182, 197, 198, 215, 263, 264], the proceedings, [5, 41, 73, 91, 117, 212, 237, 259], and survey articles [88, 148, 152, 155, 214] and refer to the references therein. This book would not have been possible without the direct or indirect contributions of numerous colleagues and students. But, first of all, I would like to thank my father for his ability to stimulate my interest and love of mathematics over the years. Also, I am deeply indebted to my friends and teachers, Professor Dr. Rainer Kress and Professor David Colton, who introduced me to the field of scattering theory and inuenced my mathematical life in an essential way. This book is dedicated to my long friendship with them! Particular thanks are given to Dr. Frank Hettlich, Dr. Stefan Ritter, and Dipl.-Math. Markus Wartha for carefully reading the manuscript. Furthermore, I would like to thank Professor William Rundell and Dr. Martin Hanke for their manuscripts on inverse Sturm–Liouville problems and conjugate gradient methods, respectively, on which parts of Chapters 4 and 2 are based. Karlsruhe, Germany April 1996

Andreas Kirsch

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2 Regularization Theory for Equations of the First Kind 2.1 A General Regularization Theory . . . . . . . . . . . . . . . . . . 2.2 Tikhonov Regularization . . . . . . . . . . . . . . . . . . . . . . . . 2.3 Landweber Iteration . . . . . . . . . . . . . . . . . . . . . . . . . . . . 2.4 A Numerical Example . . . . . . . . . . . . . . . . . . . . . . . . . . 2.5 The Discrepancy Principle of Morozov . . . . . . . . . . . . . . 2.6 Landweber’s Iteration Method with Stopping Rule . . . . 2.7 The Conjugate Gradient Method . . . . . . . . . . . . . . . . . . 2.8 Problems . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .

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3 Regularization by Discretization . . . . . . . . . . . . . . . . . . . 3.1 Projection Methods . . . . . . . . . . . . . . . . . . . . . . . . . . . 3.2 Galerkin Methods . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 3.2.1 The Least Squares Method . . . . . . . . . . . . . . . . 3.2.2 The Dual Least Squares Method . . . . . . . . . . . . 3.2.3 The Bubnov–Galerkin Method for Coercive Operators . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 3.3 Application to Symm’s Integral Equation of the First Kind . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 3.4 Collocation Methods . . . . . . . . . . . . . . . . . . . . . . . . . . . 3.4.1 Minimum Norm Collocation . . . . . . . . . . . . . . . 3.4.2 Collocation of Symm’s Equation . . . . . . . . . . . . 3.5 Numerical Experiments for Symm’s Equation . . . . . . . 3.6 The Backus–Gilbert Method . . . . . . . . . . . . . . . . . . . . 3.7 Problems . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .

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1 Introduction and Basic Concepts 1.1 Examples of Inverse Problems . 1.2 Ill-Posed Problems . . . . . . . . . . 1.3 The Worst-Case Error . . . . . . . 1.4 Problems . . . . . . . . . . . . . . . . . .

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5 Inverse Eigenvalue Problems . . . . . . . . . . . . . . . . . . . . 5.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 5.2 Construction of a Fundamental System . . . . . . . . . . 5.3 Asymptotics of the Eigenvalues and Eigenfunctions 5.4 Some Hyperbolic Problems . . . . . . . . . . . . . . . . . . . . 5.5 The Inverse Problem . . . . . . . . . . . . . . . . . . . . . . . . 5.6 A Parameter Identification Problem . . . . . . . . . . . . 5.7 Numerical Reconstruction Techniques . . . . . . . . . . . 5.8 Problems . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .

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6 An Inverse Problem in Electrical Impedance Tomography . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 6.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 6.2 The Direct Problem and the Neumann–Dirichlet Operator 6.3 The Inverse Problem . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 6.4 The Factorization Method . . . . . . . . . . . . . . . . . . . . . . . . . 6.5 Problems . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .

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4 Nonlinear Inverse Problems . . . . . . . . . . . . . . . . . . . 4.1 Local Illposedness . . . . . . . . . . . . . . . . . . . . . . . . . . 4.2 The Nonlinear Tikhonov Regularization . . . . . . . . 4.2.1 Existence of Solutions and Stability . . . . . . 4.2.2 Source Conditions And Convergence Rates 4.2.3 A Parameter-Identification Problem . . . . . . 4.2.4 A Glimpse on Extensions to Banach Spaces 4.3 The Nonlinear Landweber Iteration . . . . . . . . . . . . 4.4 Problems . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .

7 An 7.1 7.2 7.3 7.4 7.5 7.6

Inverse Scattering Problem . . . . . . . . . . . . . . . . . . . . . Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . The Direct Scattering Problem . . . . . . . . . . . . . . . . . . . . Properties of the Far Field Patterns . . . . . . . . . . . . . . . Uniqueness of the Inverse Problem . . . . . . . . . . . . . . . . . The Factorization Method . . . . . . . . . . . . . . . . . . . . . . . The Interior Transmission Eigenvalue Problem . . . . . . . 7.6.1 The Radially Symmetric Case . . . . . . . . . . . . . . . 7.6.2 Discreteness And Existence in the General Case 7.6.3 The Inverse Spectral Problem for the Radially Symmetric Case . . . . . . . . . . . . . . . . . . . . . . . . . . 7.7 Numerical Methods . . . . . . . . . . . . . . . . . . . . . . . . . . . . 7.7.1 A Simplified Newton Method . . . . . . . . . . . . . . . 7.7.2 A Modified Gradient Method . . . . . . . . . . . . . . . 7.7.3 The Dual Space Method . . . . . . . . . . . . . . . . . . . 7.8 Problems . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .

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239 239 243 256 266 272 282 284 291

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296 302 304 307 308 312

Contents

A Basic A.1 A.2 A.3 A.4 A.5 A.6 A.7 A.8 A.9 A.10

xvii

Facts from Functional Analysis . . . . . . . . . . . . . . . . . . Normed Spaces and Hilbert Spaces . . . . . . . . . . . . . . . . . . . . Orthonormal Systems . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . Linear Bounded and Compact Operators . . . . . . . . . . . . . . . . Sobolev Spaces of Periodic Functions . . . . . . . . . . . . . . . . . . Sobolev Spaces on the Unit Disc . . . . . . . . . . . . . . . . . . . . . Spectral Theory for Compact Operators in Hilbert Spaces . . The Fréchet Derivative . . . . . . . . . . . . . . . . . . . . . . . . . . . . . Convex Analysis . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . Weak Topologies . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . Problems . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .

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315 315 321 324 332 340 345 351 357 363 365

B Proofs of the Results of Section 2.7 . . . . . . . . . . . . . . . . . . . . . . . . 367 Bibliography . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 377 Index . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 397

Chapter 1

Introduction and Basic Concepts 1.1

Examples of Inverse Problems

In this section, we present some examples of pairs of problems that are inverse to each other. We start with some simple examples that are normally not even recognized as inverse problems. Most of them are taken from the survey article [152] and the monograph [111]. Example 1.1 Find a polynomial p of degree n with given zeros x1 , . . . , xn . This problem is inverse to the direct problem: Find the zeros x1 , . . . , xn of a given polynomial p. In this example, the inverse problem is easier to solve. Its solution is p(x) = c(x − x1 ) . . . (x − xn ) with an arbitrary constant c. Example 1.2 Find a polynomial p of degree n that assumes given values y1 , . . . , yn ∈ R at given points x1 , . . . , xn ∈ R. This problem is inverse to the direct problem of calculating the given polynomial at given x1 , . . . , xn . The inverse problem is the Lagrange interpolation problem. Example 1.3 Given a real symmetric n × n matrix A and n real numbers λ1 , . . . , λn , find a diagonal matrix D such that A + D has the eigenvalues λ1 , . . . , λn . This problem is inverse to the direct problem of computing the eigenvalues of the given matrix A + D. Example 1.4 This inverse problem is used with intelligence tests: Given the first few terms a1 , a2 , . . . , ak of a sequence, find the law of formation of the sequence; that is, find an for all n! Usually, only the next two or three terms are asked for to show © Springer Nature Switzerland AG 2021 A. Kirsch, An Introduction to the Mathematical Theory of Inverse Problems, Applied Mathematical Sciences 120, https://doi.org/10.1007/978-3-030-63343-1 1

1

2

Introduction

that the law of formation has been found. The corresponding direct problem is to evaluate the sequence (an ) given the law of formation. It is clear that such inverse problems always have many solutions (from the mathematical point of view), and for this reason their use on intelligence tests has been criticized. Example 1.5 (Geological prospecting) In general, this is the problem of determining the location, shape, and/or some parameters (such as conductivity) of geological anomalies in the Earth’s interior from measurements at its surface. We consider a simple one-dimensional example and describe the following inverse problem. Determine changes ρ = ρ(x), 0 ≤ x ≤ 1, of the mass density of an anomalous region at depth h from measurements of the vertical component fv (x) of the change of force at x. ρ(x )Δx is the mass of a “volume element” at x and  (x − x )2 + h2 is its distance from the instrument. The change of gravity is described by Newton’s law of gravity f = γ rm2 with gravitational constant γ. For the vertical component, we have

Δfv (x) = γ

ρ(x )Δx h ρ(x )Δx cos θ = γ  3/2 .  2 2 (x − x ) + h (x − x )2 + h2

x 6 h

-x

θ ?

? x

0

1

-x

This yields the following integral equation for the determination of ρ: 1 fv (x) = γ h 0

ρ(x )   3/2 dx (x − x )2 + h2

for 0 ≤ x ≤ 1 .

(1.1)

We refer to [6, 105, 277] for further reading on this and related inverse problems in geological prospecting. Example 1.6 (Inverse scattering problem) Find the shape of a scattering object, given the intensity (and phase) of sound or electromagnetic waves scattered by this object. The corresponding direct problem is that of calculating the scattered wave for a given object.

1.1

Examples of Inverse Problems

3

us u

i

D

us More precisely, the direct problem can be described as follows. Let a bounded region D ⊂ RN (N = 2 or 3) be given with smooth boundary ∂D (the scattering ˆ object) and a plane incident wave ui (x) = eikθ·x , where k > 0 denotes the wave number and θˆ is a unit vector that describes the direction of the incident wave. The direct problem is to find the total field u = ui +us as the sum of the incident field ui and the scattered field us such that u + k 2 u = 0 in RN \ D ,   ∂us − ikus = O r−(N +1)/2 ∂r

u=0

on ∂D ,

for r = |x| → ∞ uniformly in

(1.2a)

x . |x|

(1.2b)

For acoustic scattering problems, v(x, t) = u(x)e−iωt describes the pressure and k = ω/c is the wave number with speed of sound c. For suitably polarized time harmonic electromagnetic scattering problems, Maxwell’s equations reduce to the two-dimensional Helmholtz equation Δu + k 2 u = 0 for the components of the electric (or magnetic) field u. The wave number k is given in terms of the √ dielectric constant ε and permeability μ by k = εμ ω. In both cases, the radiation condition (1.2b) yields the following asymptotic behavior: us (x) =

  exp(ik|x|) u∞ (ˆ x) + O |x|−(N +1)/2 (N −1)/2 |x|

as |x| → ∞ ,

where x ˆ = x/|x|. The inverse problem is to determine the shape of D when the far field pattern u∞ (ˆ x) is measured for all x ˆ on the unit sphere in RN . These and related inverse scattering problems have various applications in computer tomography, seismic and electromagnetic exploration in geophysics, and nondestructive testing of materials, for example. An inverse scattering problem of this type is treated in detail in Chapter 7. Standard literature on these direct and inverse scattering problems are the monographs [53, 55, 176] and the survey articles [50, 248]. Example 1.7 (Computer tomography) The most spectacular application of the Radon transform is in medical imaging. For example, consider a fixed plane through a human body. Let ρ(x, y) denote the change of density at the point (x, y), and let L be any line in the plane. Suppose that we direct a thin beam of X–rays into the body along L and measure how much of the intensity is attenuated by going through the body.

4

Introduction

L

y

s

δ

x

Let L be parametrized by (s, δ), where s ∈ R and δ ∈ [0, π). The ray Ls,δ has the coordinates seiδ + iueiδ ∈ C, u ∈ R , where we have identified C with R2 . The attenuation of the intensity I is approximately described by dI = −γρI du with some constant γ. Integration along the ray yields u   ln I(u) = −γ ρ seiδ + iteiδ dt u0

or, assuming that ρ is of compact support, the relative intensity loss is given by ∞   ln I (∞) = −γ ρ seiδ + iteiδ dt . −∞

In principle, from the attenuation factors we can compute all line integrals ∞   (Rρ)(s, δ) := (1.3) ρ seiδ + iueiδ du, s ∈ R, δ ∈ [0, π) . −∞

Rρ is called the Radon transform of ρ. The direct problem is to compute the Radon transform Rρ when ρ is given. The inverse problem is to determine the density ρ for a given Radon transform Rρ (that is, measurements of all line integrals). The problem simplifies in the following special case, where we assume that ρ is  radially symmetric and we choose only vertical rays. Then ρ = ρ(r), r = x2 + y 2 , and the ray Lx passing through (x, 0) can be parametrized by (x, u), u ∈ R. This leads to (the factor 2 is due to symmetry) ∞  x2 + u2 du . V (x) := ln I(∞) = −2γ ρ 0

Again, we assume √ that ρ is of compact support in {x : |x| ≤ R}. The change of variables u = r2 − x2 leads to ∞ V (x) = −2γ x

r √ ρ(r) dr = −2γ 2 r − x2

R √ x

r2

r ρ(r) dr . − x2

(1.4)

1.1

Examples of Inverse Problems

5

A further change of variables z = R2 − r2 and y = R2 − x2 transforms √ this equation into the following Abel’s integral equation for the function z → ρ( R2 − z): y √ 2  ρ( R − z) 2 √ dz, V ( R − y) = −γ y−z

0 ≤ y ≤ R.

(1.5)

0

The standard mathematical literature on the Radon transform and its applications are the monographs [128, 130, 206]. We refer also to the survey articles [131, 183, 185, 192]. The following example is due to Abel himself. Example 1.8 (Abel’s integral equation) Let a mass element move along a curve Γ from a point p1 on level h > 0 to a point p0 on level h = 0. The only force acting on this mass element is the gravitational force mg.

y 6

r p1

h r p

r p0

Γ: x = ψ(y)

-x

The direct problem is to determine the time T in which the element moves from p1 to p0 when the curve Γ is given. In the inverse problem, one measures the time T = T (h) for several values of h and tries to determine the curve Γ. Let the curve be parametrized by x = ψ(y). Let p have the coordinates (ψ(y), y). By conservation of energy, that is, m 2 v + m g y = const = m g h , E+U = 2 we conclude for the velocity that  ds = v = 2g(h − y) . dt The total time T from p1 to p0 is p1 h

ds 1 + ψ  (y)2 = dy for h > 0 . T = T (h) = v 2g (h − y) p0

0

6

Introduction

 √ Set φ(y) = 1 + ψ  (y)2 and let f (h) := T (h) 2g be known (measured). Then we have to determine the unknown function φ from Abel’s integral equation h 0

φ(y) √ dy = f (h) h−y

for h > 0 .

(1.6)

A similar—but more important—problem occurs in seismology. One studies the problem to determine the velocity distribution c of the Earth from measurements of the travel times of seismic waves (see [29]). For further examples of inverse problems leading to Abel’s integral equations, we refer to the lecture notes by R. Gorenflo and S. Vessella [108], the monograph [198], and the papers [179, 270]. Example 1.9 (Backwards heat equation) Consider the one-dimensional heat equation ∂ 2 u(x, t) ∂u(x, t) = , ∂t ∂x2

(x, t) ∈ (0, π) × R>0 ,

(1.7a)

t ≥ 0,

(1.7b)

0 ≤ x ≤ π.

(1.7c)

with boundary conditions u(0, t) = u(π, t) = 0, and initial condition u(x, 0) = u0 (x),

The separation of variables leads to the (formal) solution u(x, t) =

∞ 

−n2 t

an e

sin(nx)

with

n=1

2 an = π

π u0 (y) sin(ny)dy .

(1.8)

0

The direct problem is to solve the classical initial boundary value problem: Given the initial temperature distribution u0 and the final time T , determine u(·, T ). In the inverse problem, one measures the final temperature distribution u(·, T ) and tries to determine the temperature at earlier times t < T , for example, the initial temperature u(·, 0). From the solution formula (1.8), we see that we have to determine u0 := u(·, 0) from the following integral equation: 2 u(x, T ) = π

π k(x, y) u0 (y) dy,

0 ≤ x ≤ π,

(1.9)

0

where k(x, y) :=

∞ 

2

e−n

T

sin(nx) sin(ny) .

(1.10)

n=1

We refer to the monographs [17, 175, 198] and papers [31, 43, 49, 80, 81, 94, 193, 247] for further reading on this subject.

1.1

Examples of Inverse Problems

7

Example 1.10 (Diffusion in an inhomogeneous medium) The equation of diffusion in an inhomogeneous medium (now in two dimensions) is described by the equation   ∂u(x, t) 1 = div γ(x)∇u(x, t) , ∂t c

x ∈ D, t > 0,

(1.11)

where c is a constant and γ = γ(x) is a parameter describing the medium. In the stationary case, this reduces to div (γ∇u) = 0

in D .

(1.12)

The direct problem is to solve the boundary value problem for this equation for given boundary values u|∂D and given function γ. In the inverse problem, one measures u and the flux γ∂u/∂ν on the boundary ∂D and tries to determine the unknown function γ in D. This is the problem of impedance tomography which we consider in more detail in Chapter 6. The problem of impedance tomography is an example of a parameter identification problem for a partial differential equation. Among the extensive literature on parameter identification problems, we only mention the classical papers [166, 225, 224], the monographs [15, 17, 198], and the survey article [200]. Example 1.11 (Sturm–Liouville eigenvalue problem) Let a string of length L and mass density ρ = ρ(x) > 0, 0 ≤ x ≤ L, be fixed at the endpoints x = 0 and x = L. Plucking the string produces tones due to vibrations. Let v(x, t), 0 ≤ x ≤ L, t > 0, be the displacement at x and time t. It satisfies the wave equation ρ(x)

∂ 2 v(x, t) ∂ 2 v(x, t) = , 2 ∂t ∂x2

0 < x < L, t > 0,

(1.13)

subject to boundary conditions v(0, t) = v(L, t) = 0 for t > 0. A periodic displacement of the form   v(x, t) = w(x) a cos ωt + b sin ωt with frequency ω > 0 is called a pure tone. This form of v solves the boundary value problem (1.13) if and only if w and ω satisfy the Sturm–Liouville eigenvalue problem w (x) + ω 2 ρ(x) w(x) = 0 , 0 < x < L ,

w(0) = w(L) = 0 .

(1.14)

The direct problem is to compute the eigenfrequencies ω and the corresponding eigenfunctions for known function ρ. In the inverse problem, one tries to determine the mass density ρ from a number of measured frequencies ω. We see in Chapter 5 that parameter estimation problems for parabolic and hyperbolic initial boundary value problems are closely related to inverse spectral problems.

8

Introduction

Example 1.12 (Inverse Stefan problem) The physicist Stefan (see [253]) modeled the melting of arctic ice in the summer by a simple one-dimensional model. In particular, consider a homogeneous block of ice filling the region x ≥  at time t = 0. The ice starts to melt by heating the block at the left end. Thus, at time t > 0 the region between x = 0 and x = s(t) for some s(t) > 0 is filled with water, and the region x ≥ s(t) is filled with ice.

t T 6

x = s(t)

water

ice 

- x

Let u(x, t) be the temperature at 0 < x < s(t) and time t. Then u satisfies the one-dimensional heat equation ∂u(x, t) ∂ 2 u(x, t) = ∂t ∂x2

in D := {(x, t) ∈ R2 : 0 < x < s(t), t > 0}

(1.15)

∂ subject to boundary conditions ∂x u(0, t) = f (t) and u(s(t), t) = 0 for t ∈ [0, T ] and initial condition u(x, 0) = u0 (x) for 0 ≤ x ≤ . Here, u0 describes the initial temperature and f (t) the heat flux at the left boundary x = 0. The speed at which the interface between water and ice moves is proportional to the heat flux. This is described by the following Stefan condition: ∂u(s(t), t) ds(t) = − for t ∈ [0, T ] . (1.16) dt ∂x The direct problem is to compute the curve s when the boundary data f and u0 are given. In the inverse problem, one has given a desired curve s and tries to reconstruct u and f (or u0 ). We refer to the monographs [39, 198] and the classical papers [40, 95] for a detailed introduction to Stefan problems.

In all of these examples, we can formulate the direct problem as the evaluation of an operator K acting on a known “model” x in a model space X and the inverse problem as the solution of the equation K(x) = y: Direct problem: Inverse problem:

given x (and K), evaluate K(x). given y (and K), solve K(x) = y for x.

In order to formulate an inverse problem, the definition of the operator K, including its domain and range, has to be given. The formulation as an operator equation allows us to distinguish among finite, semifinite, and infinitedimensional, linear and nonlinear problems.

1.2

Ill-Posed Problems

9

In general, the evaluation of K(x) means solving a boundary value problem for a differential equation or evaluating an integral. For more general and “philosophical” aspects of inverse theory, we refer to [7, 214].

1.2

Ill-Posed Problems

For all of the pairs of problems presented in the last section, there is a fundamental difference between the direct and the inverse problems. In all cases, the inverse problem is ill-posed or improperly posed in the sense of Hadamard, while the direct problem is well-posed. In his lectures published in [115], Hadamard claims that a mathematical model for a physical problem (he was thinking in terms of a boundary value problem for a partial differential equation) has to be properly posed or well-posed in the sense that it has the following three properties: 1. There exists a solution of the problem (existence). 2. There is at most one solution of the problem (uniqueness). 3. The solution depends continuously on the data (stability). Mathematically, the existence of a solution can be enforced by enlarging the solution space. The concept of distributional solutions of a differential equation is an example. If a problem has more than one solution, then information about the model is missing. In this case, additional properties, such as sign conditions, can be built into the model. The requirement of stability is the most important one. If a problem lacks the property of stability, then its solution is practically impossible to compute because any measurement or numerical computation is polluted by unavoidable errors: thus the data of a problem are always perturbed by noise! If the solution of a problem does not depend continuously on the data, then in general the computed solution has nothing to do with the true solution. Indeed, there is no way to overcome this difficulty unless additional information about the solution is available. Here, we remind the reader of the following statement (see Lanczos [171]): A lack of information cannot be remedied by any mathematical trickery! Mathematically, we formulate the notation of well-posedness in the following way. Definition 1.13 (Well-posedness) Let X and Y be normed spaces, and K : X → Y a linear operator. The equation Kx = y is called properly posed or well-posed if the following holds: 1. Existence: For every y ∈ Y , there is (at least one) x ∈ X such that Kx = y. 2. Uniqueness: For every y ∈ Y , there is at most one x ∈ X with Kx = y.

10

Introduction 3. Stability: The solution x depends continuously on y; that is, for every sequence (xn ) in X with Kxn → Kx (n → ∞), it follows that xn → x (n → ∞).

Equations for which (at least) one of these properties does not hold are called improperly posed or ill-posed. In Chapter 4, we will extend this definition to local ill-posedness of nonlinear problems. It is important to specify the full triple (X, Y, K) and their norms. Existence and uniqueness depend only on the algebraic nature of the spaces and the operator, that is, whether the operator is onto or one-to-one. Stability, however, depends also on the topologies of the spaces, that is, whether the inverse operator K −1 : Y → X is continuous. These requirements are not independent of each other. For example, due to the open mapping theorem (see Theorem A.27 of Appendix A.3), the inverse operator K −1 is automatically continuous if K is linear and continuous and X and Y are Banach spaces. As an example for an ill-posed problem, we study the classical example given by Hadamard in his famous paper [115]. Example 1.14 (Cauchy’s problem for the Laplace equation) Find a solution u of the Laplace equation Δu(x, y) :=

∂ 2 u(x, y) ∂ 2 u(x, y) + = 0 in R × [0, ∞) ∂x2 ∂y 2

(1.17a)

that satisfies the “initial conditions” u(x, 0) = f (x),

∂ u(x, 0) = g(x) , ∂y

x ∈ R,

(1.17b)

where f and g are given functions. Obviously, the (unique) solution for f (x) = 0 and g(x) = n1 sin(nx) is given by u(x, y) =

1 sin(nx) sinh(ny), n2

x ∈ R, y ≥ 0 .

Therefore, we have sup {|f (x)| + |g(x)|} = x∈R

but sup |u(x, y)| = x∈R

1 −→ 0, n

1 sinh(ny) −→ ∞, n2

n → ∞,

n→∞

for all y > 0. The error in the data tends to zero while the error in the solution u tends to infinity! Therefore, the solution does not depend continuously on the data, and the problem is improperly posed.

1.2

Ill-Posed Problems

11

Many inverse problems and some of the examples of the last section (for further examples, we refer to [111]) lead to integral equations of the first kind with continuous or weakly singular kernels. Such integral operators are compact with respect to any reasonable topology. The following example will often serve as a model case in these lectures. Example 1.15 (Differentiation) The direct problem is to find the antiderivative y with y(0) = 0 of a given continuous function x on [0, 1], that is, compute t x(s) ds ,

y(t) =

t ∈ [0, 1] .

(1.18)

0

In the inverse problem, we are given a continuously differentiable function y on [0, 1] with y(0) = 0 and want to determine x = y  . This means we have to solve the integral equation Kx = y, where K : C[0, 1] → C[0, 1] is defined by t x(s) ds,

(Kx)(t) :=

t ∈ [0, 1] ,

for x ∈ C[0, 1] .

(1.19)

0

Here, we equip C[0, 1] with the supremum norm x ∞ := max |x(t)|. The 0≤t≤1

solution of Kx = y is just the derivative x = y  , provided y(0) = 0 and y is continuously differentiable! If x is the exact solution of Kx = y, and if we perturb y in the norm · ∞ , then the perturbed right-hand side y˜ doesn’t have to be differentiable, and even if it is the solution of the perturbed problem is not necessarily close to the exact solution. We can, for example, perturb y by δ sin(t/δ 2 ) for small δ. Then the error of the data (with  respect to · ∞ ) is δ and the error in the solution is 1/δ. The problem K, C[0, 1], C[0, 1] is therefore ill-posed.

Now we choose a different space Y := y ∈ C 1 [0, 1] : y(0) = 0 for the righthand side and equip Y with the stronger norm x C 1 := max |x (t)|. If the 0≤t≤1

right-hand side  is perturbed with respect to this norm · C 1 , then the problem K, C[0, 1], Y is well-posed because K : C[0, 1] → Y is boundedly invertible. This example again illustrates the fact that well-posedness depends on the topology. In the numerical treatment of integral equations, a discretization error cannot be avoided. For integral equations of the first kind, a “naive” discretization usually leads to disastrous results as the following simple example shows (see also [267]). Example 1.16 The integral equation 1 0

ets x(s) ds = y(t) ,

0 ≤ t ≤ 1,

(1.20)

12

Introduction

with y(t) = (exp(t + 1) − 1)/(t + 1), is uniquely solvable by x(t) = exp(t). We approximate the integral by the trapezoidal rule ⎛ ⎞ 1 n−1  1 1 ets x(s) ds ≈ h ⎝ x(0) + et x(1) + ejht x(jh)⎠ 2 2 j=1 0

with h := 1/n. For t = ih, we obtain the linear system ⎛ ⎞ n−1  2 1 1 h ⎝ x0 + eih xn + ejih xj ⎠ = y(ih) , i = 0, . . . , n . 2 2 j=1

(1.21)

Then xi should be an approximation to x(ih). The following table lists the error between the exact solution x(t) and the approximate solution xi for t = 0, 0.25, 0.5, 0.75, and 1. Here, i is chosen such that ih = t. t 0 0.25 0.5 0.75 1

n=4 0.44 0.67 0.95 1.02 1.09

n=8 0.47 2.03 4.74 3.08 1.23

n = 16 1.30 39.02 15.34 15.78 0.91

n = 32 41.79 78.39 1.72 2.01 20.95

We see that the approximations have nothing to do with the true solution and become even worse for finer discretization schemes. In the previous two examples, the problem was to solve integral equations of the first kind. Integral operators are compact operators in many natural topologies under very weak conditions on the kernels. The next theorem implies that linear equations of the form Kx = y with compact operators K are always ill-posed. Theorem 1.17 Let X, Y be normed spaces and K : X → Y be a linear compact operator with nullspace N (K) := {x ∈ X : Kx = 0}. Let the dimension of the factor space X/N (K) be infinite. Then there exists a sequence (xn ) in X such that Kxn → 0 but (xn ) does not converge. We can even choose (xn ) such that xn X → ∞. In particular, if K is one-to-one, the inverse K −1 : Y ⊃ R(K) → X is unbounded. Here, R(K) := {Kx ∈ Y : x ∈ X} denotes the range of K. Proof: We set N = N (K) for abbreviation. The factor space X/N is a

normed space with norm [x] := inf x + z X : z ∈ N since the nullspace ˜ : X/N → Y , defined by K([x]) ˜ is closed. The induced operator K := Kx, ˜ −1 : Y ⊃ [x] ∈ X/N , is well-defined, compact, and one-to-one. The inverse K ˜ : X/N → ˜ −1 K R(K) → X/N is unbounded since otherwise the identity I = K

1.3

The Worst-Case Error

13

X/N would be compact as a composition of a bounded and a compact operator (see Theorem A.34). This would contradict the assumption that the dimension ˜ −1 is unbounded, there of X/N is infinite (see again Theorem A.34). Because K exists a sequence ([zn ]) in X/N with Kzn → 0 and [zn ] =1. We choose vn ∈ N such that zn + vn X ≥ 12 and set xn := (zn + vn )/ Kzn . Then  Kxn → 0 and xn X → ∞.

1.3

The Worst-Case Error

We come back to Example 1.15 of the previous section: Determine x ∈ C[0, 1] t such that 0 x(s) ds = y(t) for all t ∈ [0, 1]. An obvious question is: How large could the error be in the worst case if the error in the right side y is at most δ? The answer is already given by Theorem 1.17: If the errors are measured in norms such that the integral operator is compact, then the solution error could be arbitrarily large. For the special Example 1.15, we have constructed explicit perturbations with this property. However, the situation is different if additional information is available. Before we study the general case, we illustrate this observation for a model example. Let y and y˜ be twice continuously differentiable and let a number E > 0 be available with y  ∞ ≤ E . (1.22) y  ∞ ≤ E and ˜ Set z := y˜ − y and assume that z  (0) = z(0) = 0 and z  (t) ≥ 0 for t ∈ [0, 1]. Then we estimate the error x ˜ − x in the solution of Example 1.15 by 2  x ˜(t) − x(t)



2

= z (t)

t = 0

t ≤

4E

d   2 z (s) ds = 2 ds

t

z  (s) z  (s) ds

0

z  (s) ds = 4 E z(t).

0

Therefore, under √ the above assumptions on z = y˜ − y we have shown that ˜ x − x ∞ ≤ y − y ∞ ≤ δ and E is a bound as in (1.22). In this √ 2 E δ if ˜ example, 2 E δ is a bound on the worst-case error for an error δ in the data and the additional information x ∞ = y  ∞ ≤ E on the solution. We define the following quite generally. Definition 1.18 Let K : X → Y be a linear bounded operator between Banach ˆ that is, there ˆ ⊂ X a subspace, and · ˆ a “stronger” norm on X; spaces, X X ˆ Then we define exists c > 0 such that x X ≤ c x Xˆ for all x ∈ X.     ˆ Kx Y ≤ δ, x ˆ ≤ E , F δ, E, · Xˆ := sup x X : x ∈ X, (1.23) X   and call F δ, E, · Xˆ the worst-case error for the error δ in the data and a priori information x Xˆ ≤ E.

14

Introduction

  ˆ F δ, E, · Xˆ depends on the operator K and the norms in X, Y , and X. It is desirable that this worst-case error not only converges to zero as δ tends to zero but that it is of order δ. This is certainly true (even without a priori information) for boundedly invertible operators, as is readily seen from the inequality x X ≤ K −1 L(Y,X) Kx Y . For compact operators K, however, and norm · Xˆ = · X , this worst-case error does not converge (see the following lemma), and one is forced to take a stronger norm · Xˆ . Lemma 1.19 Let K : X → Y be linear and compact and assume that X/N (K) is infinite-dimensional.  Then for every E > 0, there exists c > 0 and δ0 > 0  such that F δ, E, · X ≥ c for all δ ∈ (0, δ0 ).  Proof: Assume that there exists a sequence δn → 0 such that F δn , E, ·  ˜ : X/N (K) → Y be again the induced operator in X → 0 as n → ∞. Let K   ˜ −1 is bounded: Let K ˜ [xm ] = Kxm → 0. the factor space. We show that K  Then there exists a subsequence xmn with Kxmn Y ≤ δn for all n. We set  zn :=

xmn ,

if xmn X ≤ E,

E xmn −1 X xmn ,

if xmn X > E.

Then zn X ≤ E and Kzn Y ≤ δn for all n. Because the worst-case error tends to zero, we also conclude that zn X → 0. From this, we see that zn = xmn for sufficiently large n; that is, xmn → 0 as n → ∞. This argument, applied to every subsequence of the original sequence (xm ), yields that xm tends to zero ˜ −1 is bounded on the range R(K) of K. This, however, for m → ∞; that is, K contradicts the assertion of Theorem 1.17.  In the following analysis, we make use of the singular value decomposition of the operator K (see Appendix A.6, Definition A.56). Therefore, we assume from now on that X and Y are Hilbert spaces. In many applications X and Y are Sobolev spaces; that is, spaces of measurable functions such that their (generalized) derivatives are square integrable. Sobolev spaces of functions of one variable can be characterized as follows: ⎧ ⎫ t ⎨ ⎬ H p (a, b) := x ∈ C p−1 [a, b] : x(p−1) (t) = α + ψ ds, α ∈ R, ψ ∈ L2 ⎩ ⎭ a

for p ∈ N. Example 1.20 (Differentiation) As an example, we study differentiation and set X = Y = L2 (0, 1), t (Kx)(t) =

x(s) ds , 0

t ∈ (0, 1), x ∈ L2 (0, 1) ,

(1.24)

1.3

The Worst-Case Error

15

and ˆ1 X ˆ2 X

{x ∈ H 1 (0, 1) : x(1) = 0}, {x ∈ H 2 (0, 1) : x(1) = 0, x (0) = 0}.

:= :=

(1.25a) (1.25b)

ˆ 1 , and x 2 := x L2 for x ∈ X ˆ 1 . Then We define x 1 := x L2 for x ∈ X the norms · j , j = 1, 2, are stronger than · L2 (see Problem 1.2), and we can prove for every E > 0 and δ > 0: √     F δ, E, · 1 ≤ δ E and F δ, E, · 2 ≤ δ 2/3 E 1/3 . (1.26) From this result, we observe that the possibility to reconstruct x is dependent on the smoothness of the solution. We come back to this remark in a more general setting (Theorem 1.21). We will also see that these estimates are asymptotically sharp; that is, the exponent of δ cannot be increased. Proof of (1.26): First, assume that x ∈ H 1 (0, 1) with x(1) = 0. Partial integration, which is easily seen to be allowed for H 1 -functions and the Cauchy– Schwarz inequality, yields x 2L2

1 =

x(t) x(t) dt 0

1 = −

⎡ x (t) ⎣

0

⎤

⎡

x(s) ds⎦ dt + ⎣x(t)

0

1 = −

t

t

⎤t=1 x(s) ds⎦

0

t=0

x (t) (Kx)(t) dt ≤ Kx L2 x L2 .

(1.27)

0

This yields the first estimate. Now let x ∈ H 2 (0, 1) such that x(1) = 0 and x (0) = 0. Using partial integration again, we estimate x 2L2

1 =

x (t) x (t) dt

0

1 = −

 t=1 x(t) x (t) dt + x(t) x (t) t=0

0

1 = −

x(t) x (t) dt ≤ x L2 x L2 .

0

Now we substitute this into the right-hand side of (1.27):   x 2L2 ≤ Kx L2 x L2 ≤ Kx L2 x L2 x L2 . From this, the second estimate of (1.26) follows.



16

Introduction

This example is typical in the sense that integral operators are often smoothing. We can define an abstract “smoothness” of an element x ∈ X with respect to a compact operator K : X → Y by requiring that x ∈ R(K ∗ ) or x ∈ R(K ∗ K)  ∗ σ/2 for some real σ > 0. The operator or, more generally, x ∈ R (K K) (K ∗ K)σ/2 from X into itself is defined as  (K ∗ K)σ/2 x = μσj (x, xj )X xj , x ∈ X , j∈J

A.6, where {μj , xj , yj : j ∈ J} is a singular system for K (see Appendix   Theorem A.57 and formula (A.47)). We note that R(K ∗ ) = R (K ∗ K)1/2 .  ∗ σ/2  Picard’s Theorem (Theorem A.58) yields that x ∈ R (K K) is equivalent ! |(x,xj )X |2 to < ∞ which is indeed a smoothness assumption in concrete j μ2σ j applications. We refer to Example A.59 and the definition of Sobolev spaces of periodic functions (Section A.4) where smoothness is expressed by a decay of the Fourier coefficients. Theorem 1.21 Let X and Y be Hilbert spaces, and K : X → Y linear, compact, and one-to-one with dense range R(K). Let K ∗ : Y → X be the adjoint operator. "  " ˆ 1 . Then ˆ 1 := R(K ∗ ) and x 1 := " K ∗ −1 x" for x ∈ X (a) Set X Y √   F δ, E, · 1 ≤ δ E .

Furthermore, every E > 0 there exists a sequence δj → 0 such that for   F δj , E, · 1 = δj E; that is, this estimate is asymptotically sharp. " "  ˆ 2 := R(K ∗ K) and x 2 := " K ∗ K −1 x" for x ∈ X ˆ 2 . Then (b) Set X X   F δ, E, · 2 ≤ δ 2/3 E 1/3 ,

 and for every E > 0 there exists a sequence δj → 0 such that F δj , E, ·  2/3 2 = δj E 1/3 .   ˆ σ := R (K ∗ K)σ/2 and x σ := (c) More generally, for some σ > 0 define X " ∗ −σ/2 " ˆ σ . Then " K K x" for x ∈ X X

  F δ, E, · σ ≤ δ σ/(σ+1) E 1/(σ+1) ,

 and for every E > 0 there exists a sequence δj → 0 such that F δj , E, ·  σ/(σ+1) 1/(σ+1) E . σ = δj The norms · 1 , · 2 , and · σ are well-defined because K ∗ and (K ∗ K)σ/2 are one-to-one. In concrete examples, the assumptions x ∈ R(K ∗ ) and x ∈ R((K ∗ K)σ/2 are smoothness assumptions on the exact solution x (together

1.3

The Worst-Case Error

17

with boundary conditions) as we have mentioned already before. In the pret ceding example, where (Kx)(t) = 0 x(s) ds, the spaces R(K ∗ ) and R(K ∗ K) ˆ 1 and X ˆ 2 defined in (1.25a) and (1.25b) (see coincide with the Sobolev spaces X Problem 1.3). ˆ 1 with Kx Y ≤ δ and x 1 ≤ E; Proof of Theorem 1.21: (a) Let x = K ∗ z ∈ X that is, z Y ≤ E. Then x 2X = (K ∗ z, x)X = (z, Kx)Y ≤ z Y Kx Y ≤ E δ . This proves the first estimate. Now let {μj , xj , yj : j ∈ J} be a singular system xj and for K (see Appendix A.6, Theorem A.57). Set x ˆj = E K ∗ yj = μj E  2 xj 1 = E, K x ˆj Y = δj , and ˆ xj X = μj E = δj E. δj := μj E → 0. Then ˆ This proves part (a). Part (b) is proven similarly or as a special case (σ = 2) of part (c). ! (c) With a singular system {μj , xj , yj : j ∈ J}, we have x 2X = j∈J |ρj |2 where ρj = (x, xj )X are the expansion coefficients of x. In the following estimate, we use H¨older’s inequality with p = (σ + 1)/σ and q = σ + 1 (note that 1/p + 1/q = 1):    2/(σ+1) |ρj |2 = (|ρj | μj )2σ/(σ+1) |ρj |/μσj x 2X = j∈J

⎛ ≤ ⎝ ⎛ = ⎝

j∈J



|ρj | μj

j∈J



2pσ/(σ+1)

⎞σ/(σ+1) ⎛ |ρj |2 μ2j ⎠

⎝

j∈J

=

⎞1/p ⎛ ⎞1/q   2q/(σ+1) ⎠ ⎝ ⎠ |ρj |/μσj j∈J



⎞1/(σ+1)

⎠ |ρj |2 μ−2σ j

j∈J

2σ/(σ+1) 2/(σ+1) Kx Y (K ∗ K)−σ/2 X .

This ends the proof.



Next, we consider Example 1.9 again. We are given the parabolic initial boundary value problem ∂ 2 u(x, t) ∂u(x, t) = , 0 < x < π, t > 0 , ∂t ∂x2 u(0, t) = u(π, t) = 0, t > 0, u(x, 0) = u0 (x), 0 < x < π . In the inverse problem, we know the final temperature distribution u(x, T ), 0 ≤ x ≤ π, and we want to determine the temperature u(x, τ ) at time τ ∈ (0, T ). As additional information, we also assume the knowledge of E > 0 with u(·, 0) L2 ≤ E. The solution of the initial boundary value problem is given by the series π ∞ 2  −n2 t e sin(nx) u0 (y) sin(ny) dy, u(x, t) = π n=1 0

0 ≤ x ≤ π, t > 0 .

18

Introduction

We denote the unknown function by v := u(·, τ ), set X = Y = L2 (0, π), and  # ∞ ∞   2 2 −n τ 2 ˆ := v ∈ L (0, π) : v = X an e sin(n·) with a 1 with 1/p + 1/q = 1 to be specified in a moment): ∞ 2 π |an |2 e−2n τ 2 n=1

=

≤

∞  2 π |an |2/q |an |2/p e−2n τ 2 n=1 &1/q % ∞ &1/p % ∞  2 π π |an |2 |an |2 e−2pn τ . 2 n=1 2 n=1

We now choose p = T /τ . Then 1/p = τ /T and 1/q = 1 − τ /T . This yields the assertion. 

1.3

The Worst-Case Error

19

The next chapter is devoted to the construction of regularization schemes ˆ that are asymptotically optimal in the sense that, under the information x ∈ X, " " " " x Xˆ ≤ E, and y˜ − y" Y ≤ "δ, an approximation x ˜ and a constant c > 0 are ˜ − x"X ≤ c F(δ, E, · Xˆ ). constructed such that "x As the first tutorial example, we consider the problem of numerical differentiation; see Examples 1.15 and 1.20. Example 1.22 t Let again, as in Examples 1.15 and 1.20, (Kx)(t) = 0 x(s)ds, t ∈ (0, 1). Solving Kx = y is equivalent to differentiating y. We fix h ∈ (0, 1/2) and define the one-sided difference quotient by  1  h y(t + h) − y(t) , 0 < t < 1/2, v(t) =   1 h y(t) − y(t − h) , 1/2 < t < 1, for any y ∈ L2 (0, 1). First, we estimate v − y  L2 for smooth functions y; that is, y ∈ H 2 (0, 1). From Taylor’s formula (see Problem 1.4), we have y(t ± h) = y(t) ± y  (t) h +

t±h  (t ± h − s) y  (s) ds ; t

that is, 

v(t) − y (t)

=

=

1 h 1 h

t+h  (t + h − s) y  (s) ds t

h

τ y  (t + h − τ ) dτ

0

for t ∈ (0, 1/2) and analogously for t ∈ (1/2, 1). Hence, we estimate 1/2 h |v(t) − y  (t)|2 dt 2

0

h h =

≤

≤

0

0

0

0

⎡

⎤ 1/2 ⎢ ⎥ τ s ⎣ y  (t + h − τ ) y  (t + h − s) dt⎦ dτ ds 0

) ) * * 1/2 h h * 1/2 * * *  2 + τs |y (t + h − τ )| dt+ |y  (t + h − s)|2 dt dτ ds 0

0

⎡ h ⎤2  1 4  2 h y L2 , y  2L2 ⎣ τ dτ ⎦ = 4 0

20

Introduction

and analogously for h2

1 1/2

|v(t) − y  (t)|2 dt. Summing these estimates yields 1 v − y  L2 ≤ √ E h , 2

where E is some bound on y  L2 . Now we treat the situation with errors. Instead of y(t) and y(t ± h), we measure y˜(t) and y˜(t ± h), respectively. We assume that ˜ y − y L2 ≤ δ. Instead   of v(t), we compute v˜(t) = ± y˜(t ± h) − y˜(t) /h for t ∈ (0, 1/2) or t ∈ (1/2, 1), respectively. Because |˜ v (t) − v(t)| ≤

|˜ y (t) − y(t)| |˜ y (t ± h) − y(t ± h)| + , h h

we conclude that ˜ v − v L2 ≤ 2δ/h. Therefore, the total error due to the error on the right-hand side and the discretization error is v − v L2 + v − y  L2 ≤ ˜ v − y  L2 ≤ ˜

1 2δ + √ E h. h 2

(1.29)

By this estimate, it is desirable to choose the discretization parameter h as the minimum of the right-hand side of (1.29). Its minimum is obtained at $ √ √ √ h = 2 2 δ/E. This results in the optimal error ˜ v − y  L2 ≤ 2 4 2 E δ.  Summarizing, we note that the discretization parameter h should be of order δ/E if the derivative of a function is computed by the one-sided difference quotient. With this choice, the method is asymptotically optimal under the information x L2 ≤ E. The two-sided difference quotient is optimal under the a priori information x L2 ≤ E and results in an algorithm of order δ 2/3 (see Example 2.4 in the following chapter). We have carried out the preceding analysis with respect to the L2 -norm rather that the maximum norm, mainly because we present the general theory in Hilbert spaces. For this example, however, estimates with respect to · ∞ are simpler to derive (see the estimates preceding Definition 1.18 of the worst-case error). The result of this example is of practical importance: For many algorithms using numerical derivatives (for example, quasi-Newton methods in optimization), it is recommended that you choose the discretization parameter ε to be the square root of the floating-point precision of the computer because a onesided difference quotient is used.

1.4

Problems

1.1 Show that equations (1.1) and (1.20) have at most one solution. Hints: Extend ρ in (1.1) by zero into R, and apply the Fourier transform and the convolution theorem. For (1.20) use results of the Laplace transform.

1.4

Problems

21

ˆ 1 and X ˆ 2 be defined by (1.25a) and (1.25b), 1.2 Let the Sobolev spaces X respectively. Define the bilinear forms by 1 (x, y)1 :=



1



x (t) y (t) dt

and

(x, y)2 :=

0

x (t) y  (t) dt

0

ˆ 1 and X ˆ 2 , respectively. Prove that X ˆ j are Hilbert spaces with respect on X ˆj , to the inner products (·, ·)j , j = 1, 2, and that x L2 ≤ x j for all x ∈ X j = 1, 2. 1.3 Let K : L2 (0, 1) → L2 (0, 1) be defined by (1.19). Show that the ranges ˆ 1 and X ˆ 2 defined by R(K ∗ ) and R(K ∗ K) coincide with the spaces X (1.25a) and (1.25b), respectively. 1.4 Prove the following version of Taylor’s formula by induction with respect to n and partial integration: Let y ∈ H n+1 (a, b) and t, t + h ∈ [a, b]. Then y(t + h) =

n  y (k) (t) k=0

k!

hk + Rn (t; h) ,

where the error term is given by 1 Rn (t; h) = n!

t+h  (t + h − s)n y (n+1) (s) ds . t

1.5 Verify the assertions of Example A.59 of Appendix A.6.

Chapter 2

Regularization Theory for Equations of the First Kind We saw in the previous chapter that many inverse problems can be formulated as operator equations of the form Kx = y , where K is a linear compact operator between Hilbert spaces X and Y over the field K = R or C. We also saw that a successful reconstruction strategy requires additional a priori information about the solution. This chapter is devoted to a systematic study of regularization strategies for solving Kx = y. In particular, we wish to investigate under which conditions they are asymptotically optimal, that is, of the same asymptotic order as the worst-case error. In Section 2.1, we introduce the general concept of regularization. In Sections 2.2 and 2.3, we study Tikhonov’s method and the Landweber iteration as two of the most important regularization strategies. In these three sections, the regularization parameter α = α(δ) is chosen a priori, that is, before we start to compute the regularized solution. We see that the optimal regularization parameter α depends on bounds of the exact solution; they are not known in advance. Therefore, it is advantageous to study strategies for the choice of α that depend on the numerical algorithm and are made during the algorithm (a posteriori). Different a posteriori choices are studied in Sections 2.5–2.7. All of them are motivated by the idea that it is certainly sufficient to compute an approximation xα,δ of the solution x such that the norm of the defect Kxα,δ − y δ is of the same order as the perturbation error δ of the right-hand side. The classical strategy, due to Morozov [194], determines α by solving a nonlinear scalar equation. To solve this equation, we still need a numerical algorithm such as the “regula falsi” or the Newton method. In Sections 2.6 and 2.7, we investigate two well-known iterative algorithms for solving linear (or nonlinear) equations: Landweber’s method (see [172]), which is the steepest

© Springer Nature Switzerland AG 2021 A. Kirsch, An Introduction to the Mathematical Theory of Inverse Problems, Applied Mathematical Sciences 120, https://doi.org/10.1007/978-3-030-63343-1 2

23

24

Regularization Theory

descent method, and the conjugate gradient method. The choices of α are made implicitly by stopping the algorithm as soon as the defect Kxm − y δ Y is less than rδ. Here, r > 1 is a given parameter. Landweber’s method and Morozov’s discrepancy principle are easy to investigate theoretically because they can be formulated as linear regularization methods. The study of the conjugate gradient method is more difficult because the choice of α depends nonlinearly on the right-hand side y. Because the proofs in Section 2.7 are very technical, we postpone them to an appendix (Appendix B).

2.1

A General Regularization Theory

For simplicity, we assume throughout this chapter that the compact operator K is one-to-one. This is not a serious restriction because we can always replace the domain X by the orthogonal complement of the kernel of K. We make the assumption that there exists a solution x∗ ∈ X of the unperturbed equation Kx∗ = y ∗ . In other words, we assume that y ∗ ∈ R(K). The injectivity of K implies that this solution is unique. In practice, the right-hand side y ∗ ∈ Y is never known exactly but only up to an error of, say, δ > 0. Therefore, we assume that we know δ > 0 and y δ ∈ Y with (2.1) y ∗ − y δ Y ≤ δ . It is our aim to “solve” the perturbed equation Kxδ = y δ .

(2.2)

In general, (2.2) is not solvable because we cannot assume that the measured data y δ are in the range R(K) of K. Therefore, the best we can hope is to determine an approximation xδ ∈ X to the exact solution x∗ that is “not much worse” than the worst-case error F(δ, E,  · Xˆ ) of Definition 1.18. An additional requirement is that the approximate solution xδ should depend continuously on the data y δ . In other words, it is our aim to construct a suitable bounded approximation R : Y → X of the (unbounded) inverse operator K −1 : R(K) → X. Definition 2.1 A regularization strategy is a family of linear and bounded operators Rα : Y −→ X, α > 0 , such that lim Rα Kx = x

α→0

for all x ∈ X ;

that is, the operators Rα K converge pointwise to the identity. From this definition and the compactness of K, we conclude the following. Theorem 2.2 Let Rα be a regularization strategy for a compact and injective operator K : X → Y where dim X = ∞. Then we have

2.1 A General Regularization Theory

25

(1) The operators Rα are not uniformly bounded; that is, there exists a sequence (αj ) with Rαj L(Y,X) → ∞ for j → ∞. (2) The sequence (Rα Kx) does not converge uniformly on bounded subsets of X; that is, there is no convergence Rα K to the identity I in the operator norm. Proof: (1) Assume, on the contrary, that there exists c > 0 such that Rα L(Y,X) ≤ c for all α > 0. From Rα y → K −1 y (α → 0) for all y ∈ R(K) and Rα yX ≤ cyY for α > 0, we conclude that K −1 yX ≤ cyY for every y ∈ R(K); that is, K −1 is bounded. This implies that I = K −1 K : X → X is compact, a contradiction to dim X = ∞. (2) Assume that Rα K → I in L(X, X). From the compactness of Rα K and Theorem A.34, we conclude that I is also compact, which again would imply that dim X < ∞.  The notation of a regularization strategy is based on unperturbed data; that is, the regularizer Rα y ∗ converges to x∗ for the exact right-hand side y ∗ = Kx∗ . Now let y ∗ ∈ R(K) be the exact right-hand side and y δ ∈ Y be the measured data with y ∗ − y δ Y ≤ δ. We define xα,δ := Rα y δ

(2.3)

as an approximation of the solution x∗ of Kx∗ = y ∗ . Then the error splits into two parts by the following obvious application of the triangle inequality: xα,δ − x∗ X

≤

Rα y δ − Rα y ∗ X + Rα y ∗ − x∗ X

≤

Rα L(Y,X) y δ − y ∗ Y + Rα Kx∗ − x∗ X

and thus xα,δ − x∗ X ≤ δRα L(Y,X) + Rα Kx∗ − x∗ X .

(2.4a)

Analogously, for the defect in the equation we have Kxα,δ − y ∗ Y ≤ δKRα L(Y ) + KRα y ∗ − y ∗ Y .

(2.4b)

These are our fundamental estimates, which we use often in the following. We observe that the error between the exact and computed solutions consists of two parts: The first term on the right-hand side of (2.4a) describes the error in the data multiplied by the “condition number” Rα L(Y,X) of the regularized problem. By Theorem 2.2, this term tends to infinity as α tends to zero. The second term denotes the approximation error (Rα − K −1 )y ∗ X at the exact right-hand side y ∗ = Kx∗ . By the definition of a regularization strategy, this term tends to zero with α. The following figure illustrates the situation (Figure 2.1).

26

Regularization Theory error

6

Rα Kx∗ − x∗ X

α∗

Rα  α

Figure 2.1: Behavior of the total error We need a strategy to choose α = α(δ) dependent on δ in order to keep the total error as small as possible. This means that we would like to minimize δ Rα L(Y,X) + Rα Kx∗ − x∗ X . The procedure is the same in every concrete situation: One has to estimate the quantities Rα L(Y,X) and Rα Kx∗ − x∗ X in terms of α and then minimize this upper bound with respect to α. Before we carry out these steps for two model examples, we introduce the following notation. Definition 2.3 A parameter choice α = α(δ) for the regularization strategy Rα is called admissible if limδ→0 α(δ) = 0 and   sup Rα(δ) y δ − xX : y δ ∈ Y, Kx − y δ Y ≤ δ → 0 , δ → 0 , for every x ∈ X. From the fundamental estimate (2.4a), we note that a parameter choice α = α(δ) is admissible if α(δ) → 0 and δRα(δ) L(Y,X) → 0 as δ → 0 and Rα Kx − xX → 0 as α → 0. Example 2.4 (Numerical differentiation by two-sided difference quotient) t Let again (Kx)(t) = 0 x(s)ds, t ∈ (0, 1); that is, solving Kx = y is equivalent to differentiating y. It is our aim to compute the derivative of y by the two-sided difference quotient (see Example 1.22 for the one-sided difference quotient). Here α = h is the stepsize, and we define ⎧ 1

 4 y t + h2 − y(t + h) − 3 y(t) , 0 < t < h2 , ⎪ h ⎪ ⎨



 1 h h (Rh y)(t) = y t + h2 − y t − h2 , h 2 < t < 1 − 2, ⎪ ⎪

 ⎩ 1 h h h 3 y(t) + y(t − h) − 4 y t − 2 , 1 − 2 < t < 1, for y ∈ L2 (0, 1). In order to prove that Rh defines a regularization strategy, it suffices to show that Rh K are uniformly bounded with respect to h in the

2.1 A General Regularization Theory

27

operator norm of L2 (0, 1) and that Rh Kx−xL2 tends to zero for smooth x (see Theorem A.29 of Appendix A.3). Later, we show convergence for x ∈ H 2 (0, 1). The fundamental theorem of calculus (or Taylor’s formula from Problem 1.4 for n = 0) yields 1 (Rh y)(t) = h

t+h/2

1 y  (s) ds = h

t−h/2

h/2

y  (s + t) ds,

−h/2

h h 0 with Rh yL2 ≤ c2 yL2 /h for all y ∈ L2 (0, 1). Estimate (2.4a) yields Rh y δ − x∗ L2 ≤ c2

δ + c1 E h2 , h

where E is a bound on (x∗ ) L2 = (y ∗ ) L2 . Minimization with respect to h of the expression on the right-hand side leads to  h(δ) = c 3 δ/E and Rh(δ) y δ − x∗ L2 ≤ c˜ E 1/3 δ 2/3 for some c > 0 and c˜ = c2 /c + c1 c2 . We observe that this strategy is asymptotically optimal for the information x L2 ≤ E because it provides an approximation xδ that is asymptotically not worse than the worst-case error (see Example 1.20). The (one- or two-sided) difference quotient uses only local portions of the function y. An alternative approach is to first smooth the function y by mollification and then to differentiate the mollified function. Example 2.5 (Numerical differentiation by mollification) t Again, we define the operator (Kx)(t) = 0 x(s) ds, t ∈ [0, 1], but now as an operator from the (closed) subspace L20 (0, 1) :=

 

1 z ∈ L2 (0, 1) : z(s) ds = 0 0

of L2 (0, 1) into L2 (0, 1). We define the Gaussian kernel ψα by 

1 √ exp −t2 /α2 , t ∈ R , α π ∞ where α > 0 denotes a parameter. Then −∞ ψα (t) dt = 1, and the convolution ψα (t) =





∞

ψα ∗ y (t) :=

∞ ψα (t − s) y(s) ds =

−∞

ψα (s) y(t − s) ds, −∞

t ∈ R,

2.1 A General Regularization Theory

29

exists and is an L2 -function for every y ∈ L2 (R). Furthermore, by Young’s inequality (see [44], p. 102), we have ψα ∗ yL2 (R) ≤ ψα L1 (R) yL2 (R) = yL2 (R)

for all y ∈ L2 (R) .

Therefore, the operators y → ψα ∗ y are uniformly bounded in L2 (R) with respect to α. We note that ψα ∗ y is infinitely often differentiable on R for every y ∈ L2 (R). We need the two convergence properties ψα ∗ z − zL2 (R) → 0 as α → 0 for every z ∈ L2 (0, 1) and ψα ∗ z − zL2 (R) ≤

√ 2 α z  L2 (0,1)

(2.5a) (2.5b)

1

for every z ∈ H (0, 1) with z(0) = z(1) = 0. Here and in the following, we identify functions z ∈ L2 (0, 1) with functions z ∈ L2 (R) where we think of them being extended by zero outside of [0, 1]. Proof of (2.5a), (2.5b): It is sufficient   to prove (2.5b) because the space z ∈ H 1 (0, 1) : z(0) = z(1) = 0 is dense in L2 (0, 1), and the operators z → ψα ∗ z are uniformly bounded from L2 (0, 1) into L2 (R). Let the Fourier transform be defined by 1 (Fz)(t) := √ 2π

∞ z(s) eist ds,

t ∈ R,

−∞

for z ∈ S, where the Schwarz space S is defined by    p (q)  ∞   S := z ∈ C (R) : sup t z (t) < ∞ for all p, q ∈ N0 . t∈R

With this normalization, Plancherel’s theorem and the convolution theorem take the form (see [44]) √ FzL2 (R) = zL2 (R) , F(u ∗ z)(t) = 2π (Fu)(t) (Fz)(t), t ∈ R , for all z, u ∈ S. Because S is dense in L2 (R) the first formula allows it to extend F to a bounded operator from L2 (R) onto itself (see Theorem A.30), and both formulas hold also for z ∈ L2 (R). Now we combine these properties and conclude that    √  ψα ∗ z − zL2 (R) = F(ψα ∗ z) − Fz L2 (R) =  2π F(ψα ) − 1 Fz L2 (R) for every z ∈ L2 (0, 1). Partial integration yields that 1 F(z )(t) = √ 2π 

1



ist

z (s) e 0

it ds = − √ 2π

1 z(s) eist ds = (−it) (Fz)(t) 0

30

Regularization Theory

for all z ∈ H 1 (0, 1) with z(0) = z(1) = 0. We define the function φα by  √ 2 2 1 1  1 − 2π F(ψα ) = 1 − e−α t /4 , it it

φα (t) :=

t ∈ R.

Then we conclude that ψα ∗ z − zL2 (R)

  = φα F(z  )L2 (R) ≤ φα ∞ F(z  )L2 (R) = φα ∞ z  L2 (0,1) .

From

    2 1 φα (t) = α 1 − e−(αt/2) 2 (αt)/2 √ and the elementary estimate 1 − exp(−τ 2 ) /τ ≤ 2 2 for all τ > 0, the desired estimate (2.5b) follows. After these preparations, we define the regularization operators Rα : L2 (0, 1) → L20 (0, 1) by (Rα y)(t)

:=

=

d (ψα ∗ y)(t) − dt (ψα ∗ y)(t) −

1

1 0

d (ψα ∗ y)(s) ds ds

(ψα ∗ y)(s) ds

0

for t ∈ (0, 1) and y ∈ L2 (0, 1). First, we note that Rα is well-defined, that is, maps L2 (0, 1) into L20 (0, 1) and is bounded. To prove that Rα is a regularization strategy, we proceed as in the previous example and show that (i) Rα yL2 ≤

4 √ α π

yL2 for all α > 0 and y ∈ L2 (0, 1),

(ii) Rα KxL2 ≤ 2 xL2 for all α > 0 and x ∈ L20 (0, 1), that is, the operators Rα K are uniformly bounded in L20 (0, 1), and √ 1 (iii) Rα Kx − xL2 ≤ 2 2 α x L2 for all α > 0 and x ∈ H00 (0, 1), where we have set 1 (0, 1) H00

 

1 1 := x ∈ H (0, 1) : x(0) = x(1) = 0, x(s) ds = 0 . 0

To prove part (i), we estimate with the Young’s inequality Rα yL2 (0,1)

≤ ≤

2 ψα ∗ yL2 (0,1) ≤ 2 ψα ∗ yL2 (R) 4 2 ψα L1 (R) yL2 (0,1) ≤ √ yL2 (0,1) α π

2.1 A General Regularization Theory

31

for all y ∈ L2 (0, 1) because ψα L1 (R)

∞ = −2

ψα (s) ds = 2 ψα (0) =

0

2 √ . α π

This proves part (i). Now let y ∈ H 1 (0, 1) with y(0) = y(1) = 0. Then, by partial integration, (ψα

1 ∗ y)(t) =

ψα (t

1 − s) y(s) ds =

0

ψα (t − s) y  (s) ds = (ψα ∗ y  )(t) .

0

Taking y = Kx, x ∈ L20 (0, 1) yields

1 (Rα Kx)(t) = (ψα ∗ x)(t) −

(ψα ∗ x)(s) ds . 0

Part (ii) now follows from Young’s inequality. Finally, we write

1 (Rα Kx)(t) − x(t) = (ψα ∗ x)(t) − x(t) −

(ψα ∗ x)(s) − x(s) ds

0

because

1 0

x(s) ds = 0. Therefore, by (2.5b),

√ Rα Kx − xL2 (0,1) ≤ 2 ψα ∗ x − xL2 (0,1) ≤ 2 2 α x L2 1 for all x ∈ H00 (0, 1). This proves part (iii). Now we conclude that Rα Kx converges to x for any x ∈ L20 (0, 1) by (ii), (iii), 1 (0, 1) in L20 (0, 1). Therefore, Rα defines a regularization and the denseness of H00 strategy. From (i) and (iii), we rewrite the fundamental estimate (2.4a) as

Rα y δ − x∗ L2 ≤

√ 4δ √ + 2 2αE α π

1 if x∗ ∈ H00 (0, 1) with (x∗ ) L2 ≤ E,y = Kx, and y δ ∈ L2 (0, 1) such that δ ∗ y − y L2 ≤ δ. The choice α = c δ/E again leads to the optimal order

√  O δE . For further applications of the mollification method, we refer to the monograph by Murio [198]. There exists an enormous number of publications on numerical differentiation. We mention only the papers [3, 69, 74, 165] and, for more general Volterra equations of the first kind, [25, 26, 76, 77, 178].

A convenient method to construct classes of admissible regularization strategies is given by filtering singular systems. Let K : X → Y be a linear

32

Regularization Theory

compact operator, and let {μj , xj , yj : j ∈ J} be a singular system for K (see Appendix A.6, Definition A.56, and Theorem A.57). As readily seen, the solution x of Kx = y is given by Picard’s theorem (see Theorem A.58 of Appendix A.6) as ∞  1 (y, yj )Y xj (2.6) x = μ j=1 j provided the series converges, that is, y ∈ R(K). This result illustrates again the influence of errors in y because the large factors 1/μj (note that μj → 0 as j → ∞) amplify the errors in the expansion coefficients (y, yj )Y . We construct regularization strategies by damping the factors 1/μj . Theorem 2.6 Let K : X → Y be compact and one-to-one with singular system {μj , xj , yj : j ∈ N} and

q : (0, ∞) × 0, KL(X,Y ) −→ R be a function with the following properties: (1) |q(α, μ)| ≤ 1 for all α > 0 and 0 < μ ≤ KL(X,Y ) . (2) For every α > 0, there exists c(α) such that |q(α, μ)| ≤ c(α) μ

for all 0 < μ ≤ KL(X,Y ) .

(3a) lim q(α, μ) = 1 for every 0 < μ ≤ KL(X,Y ) . α→0

Then the operator Rα : Y → X, α > 0, defined by Rα y :=

∞  q(α, μj )

μj

j=1

(y, yj )Y xj ,

y∈Y ,

(2.7)

is a regularization strategy with Rα L(Y,X) ≤ c(α)

and KRα L(Y ) ≤ 1. A choice α = α(δ) is admissible if α(δ) → 0 and δ c α(δ) → 0 as δ → 0. The function q is called a regularizing filter for K. Proof:

The operators Rα are bounded because we have by assumption (2) that Rα y2X

=

∞  2 1 2 q(α, μj ) 2 |(y, yj )Y | μ j j=1

≤

c(α)2

∞ 

|(y, yj )Y |2 ≤ c(α)2 y2Y ;

j=1

that is, Rα L(Y,X) ≤ c(α). From KRα y =

∞  q(α, μj ) j=1

μj

(y, yj )Y Kxj =

∞  j=1

q(α, μj ) (y, yj )Y yj ,

2.1 A General Regularization Theory

33

∞ 2 2 2 we conclude that KRα y2Y = j=1 |q(α, μj )| |(y, yj )Y | ≤ yY and thus KRα L(Y ) ≤ 1. Furthermore, from Rα Kx =

∞  q(α, μj )

μj

j=1

(Kx, yj )Y xj

and x =

∞ 

(x, xj )X xj ,

j=1

and (Kx, yj )Y = (x, K ∗ yj )X = μj (x, xj )X , we conclude that Rα Kx − x2X =

∞  2 q(α, μj ) − 1 |(x, xj )X |2 .

(2.8)

j=1

Here, K ∗ denotes the adjoint of K (see Theorem A.24). This fundamental representation will be used quite often in the following. Now let x ∈ X be arbitrary but fixed. For  > 0 there exists N ∈ N such that ∞ 

|(x, xj )X |2
0 such that [q(α, μj ) − 1]2
2. Here cσ is independent of α. It is c1 = 1/2 and c2 = 1. (b) q(α, μ) = 1−(1−a μ2 )1/α for some 0 < a < 1/K2L(X,Y ) . In this case (2) 

σ σ/2 σ/2 holds with c(α) = a/α, and (3b) is satisfied with ωσ (α) = 2a α for all σ, α > 0. (c) Let q be defined by  q(α, μ) =

1, μ2 ≥ α, 0, μ2 < α.

√ In this case (2) holds with c(α) = 1/ α, and (3b) is satisfied with ωσ (α) = ασ/2 for all σ, α > 0. Therefore, all of the functions q defined in (a), (b), and (c) are regularizing filters. Proof:

For all three cases, properties (1) and (3a) are obvious.

(a) Property (2) follows from the elementary estimate √ α, μ > 0 (which is equivalent to (μ − α)2 ≥ 0).

μ α+μ2

≤

1 √ 2 α

for all

We observe that 1 − q(α, μ) = α/(α + μ2 ). For fixed α > 0 and σ > 0, we define the function

 f (μ) = μσ 1 − q(α, μ) =

α μσ , 0 ≤ μ ≤ μ0 , α + μ2

and compute its derivative as f  (μ) = α μσ−1

(σ − 2) μ2 + ασ . (α + μ2 )2

Then f is monotonically increasing for σ ≥ 2 and thus f (μ) ≤ f (μ0 ) ≤ μσ−2 α. 0 ασ with value f (μmax ) ≤ If σ < 2, then we compute its maximum as μ2max = 2−σ cσ ασ/2 . We leave the details to the reader.

36

Regularization Theory

(b) Property (2) follows immediately from Bernoulli’s inequality:  

1/α a μ2 a μ2 1 − 1 − a μ2 , ≤ 1− 1− = α α thus |q(α, μ)| ≤



|q(α, μ)| ≤



a/α μ.

(3b) is shown by the same method as in (a) when we define

 f (μ) = μσ 1 − q(α, μ) = μσ (1 − a μ2 )1/α , 0 ≤ μ ≤ a , with derivative f  (μ) = μσ−1 (1 − a μ2 )1/α−1 σ(1 − aμ2 ) − 2aμ2 . Then aμ2max = ασ σ/2 . Again, we leave the details to the reader. 2+ασ with value f (μmax ) ≤ cσ α 2 (c) For property (2), √ it is sufficient to consider the case μ ≥ 2α. In this case, q(α,

μ) = 1 ≤μ/ α. For (3b), we consider only the case μ < α and have  μσ 1 − q(α, μ) = μσ ≤ ασ/2 .

We will see later that the regularization methods for the first two choices of q admit a characterization that avoids knowledge of the singular system. The choice (c) of q is called the spectral cutoff. The spectral cutoff solution xα,δ ∈ X is therefore defined by xα,δ =

 1 (y δ , yj )Y xj . μ j 2

μj ≥α

For this spectral cutoff solution, we combine the fundamental estimates (2.12a), (2.12b) with the previous theorem and show the following result. Theorem 2.9 (a) Let K : X → Y be a compact and injective operator with singular system {μj , xj , yj : j ∈ N}. The operators Rα y :=

 1 (y, yj )Y xj , μj 2

y∈Y ,

(2.13)

μj ≥α

√ define a regularization strategy with Rα L(Y,X) ≤ 1/ α. Thus the parameter choice α = α(δ) is admissible if α(δ) → 0 (δ → 0) and δ 2 /α(δ) → 0 (δ → 0). (b) Let Kx∗ = y ∗ and y δ ∈ Y be such that y δ − y ∗ Y ≤ δ. Furthermore, let x∗ = K ∗ z ∈ R(K ∗ ) with zY ≤ E and c > 0. For the choice α(δ) = c δ/E, we have the following error estimates for the spectral cutoff regularization.     α(δ),δ √ √ 1 ∗ x √ + c δE, (2.14a) −x X ≤ c   α(δ),δ Kx − y ∗ Y ≤ (1 + c) δ . (2.14b)

2.2 Tikhonov Regularization

37



(c) Let x∗ = (K ∗ K)σ/2 z ∈ R (K ∗ K)σ/2 for some σ > 0 with zX ≤ E. The choice α(δ) = c δ 2/(σ+1) leads to the estimates     α(δ),δ 1 ∗ σ/2 x √ +c −x X ≤ (2.14c) δ σ/(σ+1) E 1/(σ+1) , c   α(δ),δ 

Kx (2.14d) − y ∗ Y ≤ 1 + c(σ+1)/2 δ . Therefore, the spectral cutoff regularization is optimal under the information (K ∗ )−1 x∗ Y ≤ E or (K ∗ K)−σ/2 x∗ X ≤ E, respectively (if K ∗ is one-toone). Proof: Combining the fundamental estimates (2.12a), (2.12b) with Theorem 2.8 (part (c)) yields the error estimates xα,δ − x∗ X

≤

√ δ √ + α zY , α

Kxα,δ − y ∗ Y

≤

δ + αzY ,

for part (b) and xα,δ − x∗ X

≤

δ √ + ασ/2 zX , α

Kxα,δ − y ∗ Y

≤

δ + α(σ+1)/2 zX ,

for part (c). The choices α(δ) = c δ/E and α(δ) = c(δ/E)2/(σ+1) , respectively, lead to the estimates (2.14a), (2.14b) and (2.14c), (2.14d), respectively.  The general regularization concept discussed in this section can be found in many books on inverse theory [17, 110, 182]. It was not the aim of this section to study the most general theory. This concept has been extended in several directions. For example, in [84] (see also [88]) the notations of strong and weak convergence and divergence are defined, and in [182] different notations of optimality of regularization schemes are discussed . The idea of using filters has a long history [109, 265] and is very convenient for theoretical purposes. For given concrete integral operators, however, one often wants to avoid the computation of a singular system. In the next sections, we give equivalent characterizations for the first two examples without using singular systems.

2.2

Tikhonov Regularization

A common method to deal with overdetermined finite linear systems of the form Kx = y is to determine the best fit in the sense that one tries to minimize the defect Kx − yY with respect to x ∈ X for some norm in Y . If X is infinitedimensional and K is compact, this minimization problem is also ill-posed by the following lemma.

38

Regularization Theory

Lemma 2.10 Let X and Y be Hilbert spaces, K : X → Y be linear and bounded, and y ∗ ∈ Y . There exists x ˆ ∈ X with K x ˆ − y ∗ Y ≤ Kx − y ∗ Y for all ˆ = K ∗ y ∗ . Here, x ∈ X if and only if x ˆ ∈ X solves the normal equation K ∗ K x ∗ K : Y → X denotes the adjoint of K. Proof:

A simple application of the binomial theorem yields

 Kx − y ∗ 2Y − K x ˆ − y ∗ 2Y = 2 Re K x ˆ − y ∗ , K(x − x ˆ) Y + K(x − x ˆ)2Y

∗  = 2 Re K (K x ˆ − y ∗ ), x − x ˆ X + K(x − x ˆ)2Y

for all x, x ˆ ∈ X. If x ˆ satisfies K ∗ K x ˆ = K ∗ y ∗ , then Kx−y ∗ 2Y −K x ˆ−y ∗ 2Y ≥ 0, ∗ ˆ minimizes Kx−y ∗ Y , that is, x ˆ minimizes Kx−y Y . If, on the other hand, x then we substitute x = x ˆ + tz for any t > 0 and z ∈ X and arrive at 

ˆ − y ∗ ), z X + t2 Kz2Y . 0 ≤ 2t Re K ∗ (K x

 Division by t > 0 and t → 0 yields Re K ∗ (K x ˆ − y ∗ ), z X ≥ 0 for all z ∈ X; ˆ − y ∗ ) = 0, and x ˆ solves the normal equation.  that is, K ∗ (K x

As a consequence of this lemma, we should penalize the defect (in the language of optimization theory) or replace the equation of the first kind K ∗ K x ˆ= K ∗ y ∗ by an equation of the second kind (in the language of integral equation theory). Both viewpoints lead to the following minimization problem. Given the linear, bounded operator K : X → Y and y ∈ Y , determine xα ∈ X that minimizes the Tikhonov functional Jα (x) := Kx − y2Y + αx2X

for x ∈ X .

(2.15)

We prove the following theorem. Theorem 2.11 Let K : X → Y be a linear and bounded operator between Hilbert spaces and α > 0. Then the Tikhonov functional Jα has a unique minimum xα ∈ X. This minimum xα is the unique solution of the normal equation α xα + K ∗ Kxα = K ∗ y .

(2.16)

∗

The operator αI + K K is an isomorphism from X onto itself for every α > 0. Proof:

We use the following formula as in the proof of the previous lemma:

 Jα (x) − Jα (xα ) = 2 Re Kxα − y, K(x − xα ) Y + 2α Re(xα , x − xα )X + K(x − xα )2Y + αx − xα 2X 

= 2 Re K ∗ (Kxα − y) + αxα , x − xα X + K(x − xα )2Y + αx − xα 2X

(2.17)

for all x ∈ X. From this, the equivalence of the normal equation with the minimization problem for Jα is shown exactly as in the proof of Lemma 2.10. Next,

2.2 Tikhonov Regularization

39

we show that α I + K ∗ K is one-to-one for every α > 0. Let αx + K ∗ Kx = 0. Multiplication by x yields α(x, x)X + (Kx, Kx)Y = 0, that is, x = 0. Finally, we show that αI + K ∗ K is onto. Since αI + K ∗ K is one-to-one and selfadjoint, we conclude that its range is dense in X. It remains to show that the range is closed. Let zn = αxn + K ∗ Kxn converge to some z ∈ X. Then zn − zm = α(xn − xm ) + K ∗ K(xn − xm ). Multiplication of this equation by xn − xm yields αxn − xm 2X + K(xn − xm )2Y

=

(zn − zm , xn − xm )X

≤

zn − zm X xn − xm X .

From this, we conclude that αxn − xm X ≤ zn − zm X . Therefore, (xn ) is a Cauchy sequence and thus convergent to some x ∈ X which obviously satisfies  αx + K ∗ Kx = z. The solution xα of equation (2.16) can be written in the form xα = Rα y with Rα := (αI + K ∗ K)−1 K ∗ : Y −→ X . (2.18) Choosing a singular system {μj , xj , yj : j ∈ N} for the compact and injective operator K, we see that Rα y has the representation Rα y =

∞ 

∞  μj q(α, μj ) (y, y ) x = (y, yj )Y xj , j Y j 2 α + μ μj j n=0 n=0

y ∈ Y , (2.19)

with q(α, μ) = μ2 /(α + μ2 ). This function q is exactly the filter function that was studied in Theorem 2.8, part (a). Therefore, applications of Theorems 2.6 and 2.7 yield the following. Theorem 2.12 Let K : X → Y be a linear, compact, and injective operator and α > 0 and x∗ ∈ X be the exact solution of Kx∗ = y ∗ . Furthermore, let y δ ∈ Y with y δ − y ∗ Y ≤ δ. (a) The operators R √α : Y → X from (2.18) form a regularization strategy with Rα L(Y,X) ≤ 1/(2 α). It is called the Tikhonov regularization method. Rα y δ is determined as the unique solution xα,δ ∈ X of the equation of the second kind α xα,δ + K ∗ Kxα,δ = K ∗ y δ .

(2.20)

Every choice α(δ) → 0 (δ → 0) with δ 2 /α(δ) → 0 (δ → 0) is admissible. (b) Let x∗ = K ∗ z ∈ R(K ∗ ) with zY ≤ E. We choose α(δ) = c δ/E for some c > 0. Then the following estimates hold: xα(δ),δ − x∗ X

≤

Kxα(δ),δ − y ∗ Y

≤

√ √ 1 √ 1/ c + c δ E , 2 (1 + c) δ .

(2.21a) (2.21b)

40

Regularization Theory



(c) For some σ ∈ (0, 2], let x∗ = (K ∗ K)σ/2 z ∈ R (K ∗ K)σ/2 with zX ≤ E. The choice α(δ) = c (δ/E)2/(σ+1) for c > 0 leads to the error estimates   1 α(δ),δ ∗ σ/2 √ + cσ c x − x X ≤ δ σ/(σ+1) E 1/(σ+1) , (2.21c) 2 c 

Kxα(δ),δ − y ∗ Y ≤ 1 + cσ+1 c(σ+1)/2 δ . (2.21d) Here, cσ are the constants for the choice of q of part (a) of Theorem 2.8. Therefore, for σ ≤ 2 Tikhonov’s regularization method is optimal for the information (K ∗ )−1 x∗ Y ≤ E or (K ∗ K)−σ/2 x∗ X ≤ E, respectively (provided K ∗ is oneto-one). Proof: Combining the fundamental estimates (2.12a), (2.12b) with Theorem 2.8 (part (a)) yields the error estimates √ α δ √ + zY xα,δ − x∗ X ≤ 2 2 α Kxα,δ − y ∗ X

≤

δ + αzY

for part (b) and xα,δ − x∗ X

≤

δ √ + cσ ασ/2 zX 2 α

Kxα,δ − y ∗ X

≤

δ + cσ+1 α(σ+1)/2 zX

for part (c). The choices α(δ) = c δ/E and α(δ) = c(δ/E)2/(σ+1) , respectively, lead to the estimates (2.21a), (2.21b) and (2.21c), (2.21d), respectively.  From Theorem 2.12, we observe that α has to be chosen to depend on δ in such a way that it converges to zero as δ tends to zero but not as fast as δ 2 . From parts (b) and (c), we conclude that the smoother the solution x∗ is, the slower α has to tend to zero. On the other hand, the convergence can be arbitrarily slow if no a priori assumption about the solution x∗ (such as (b) or (c)) is available (see [243]). The case σ = 2 leads to the order O(δ 2/3 ) for xα(δ),δ −x∗ X . It is surprising to note that this order of convergence of Tikhonov’s regularization method 

cannot be improved even if x∗ ∈ R (K ∗ K)σ/2 for σ > 2. Indeed, we prove the following result. Theorem 2.13 Let K : X → Y be linear, compact, and one-to-one such that the range R(K) is infinite-dimensional. Furthermore, let x ∈ X, and assume that there exists a continuous function α : [0, ∞) → [0, ∞) with α(0) = 0 such that   lim xα(δ),δ − xX δ −2/3 = 0 δ→0

for every y ∈ Y with y δ − KxY ≤ δ, where xα(δ),δ ∈ X solves (2.20) for α = α(δ). Then x = 0. δ

2.2 Tikhonov Regularization

41

Proof: Assume, on the contrary, that x = 0. First, we show that α(δ) δ −2/3 → 0. Set y = Kx. From



 α(δ) I + K ∗ K xα(δ),δ − x = K ∗ (y δ − y) − α(δ) x , we estimate |α(δ)| xX ≤ KL(X,Y ) δ +



α(δ) + K2L(X,Y ) xα(δ),δ − xX .

We multiply this equation by δ −2/3 and use the assumption that xα(δ),δ tends to x faster than δ 2/3 to zero, that is, xα(δ),δ − xX δ −2/3 → 0 . This yields α(δ) δ −2/3 → 0. In the second part we construct a contradiction. Let {μj , xj , yj : j ∈ N} be a singular system for K. Define δj := μ3j

and

y δj := y + δj yj ,

j ∈ N.

Then δj → 0 as j → ∞ and, with αj := α(δj ) and xαj := (αj I + K ∗ K)−1 y, 



 xαj ,δj − x = xαj ,δj − xαj + xαj − x −1 ∗



K (δj yj ) + xαj − x = αj I + K ∗ K

 δj μj xj + xαj − x . = 2 αj + μj −2/3

Because also xαj − xX δj

→ 0, we conclude that

1/3

δj μj −→ 0, αj + μ2j

j → ∞.

But, on the other hand, 1/3

δj μj μ2j −2/3 −1 = = 1 + αj δj −→ 1, 2 2 αj + μj αj + μj This is a contradiction.

j → ∞.



This result shows that Tikhonov’s regularization method is not optimal for 

stronger “smoothness” assumptions x∗ ∈ R (K ∗ K)σ/2 for σ > 2. This is in contrast to, for example, the spectral cutoff regularization (see Theorem 2.9 above) or Landweber’s method or the conjugate gradient method, which are discussed later. The choice of α in Theorem 2.12 is made a priori, that is, before starting the computation of xα,δ by solving the least squares problem. In Sections 2.5

42

Regularization Theory

to 2.7 we study a posteriori choices of α, that is, choices of α made during the process of computing xα,δ . It is possible to choose stronger norms in the penalty term of the Tikhonov functional. Instead of (2.15), one can minimize the functional Kx − y δ 2y + αx21

on X1 ,

where  · 1 is a stronger norm (or only seminorm) on a subspace X1 ⊂ X. This was originally done by Phillips [217] and Tikhonov [261, 262] (see also [97]) for linear integral equations of the first kind. They chose the seminorm 1/2

x1 := x L2 or the H 1 -norm x1 := x2L2 + x 2L2 . By characterizing  · 1 through a singular system for K, one obtains similar convergence results as above in the stronger norm  · 1 . For further aspects of regularization with differential operators or stronger norms, we refer to [70, 119, 180, 205] and the monographs [110, 111, 182]. The interpretation of regularization by smoothing norms in terms of reproducing kernel Hilbert spaces has been observed in [133].

2.3

Landweber Iteration

Landweber [172], Fridman [98], and Bialy [18] suggested to rewrite the equation Kx = y in the form x = (I − a K ∗ K) x + a K ∗ y for some a > 0 and iterate this equation, that is, compute x0 := 0

and xm = (I − a K ∗ K) xm−1 + a K ∗ y

(2.22)

for m = 1, 2, . . . . This iteration scheme can be interpreted as the steepest descent algorithm applied to the quadratic functional x → Kx − y2Y as the following lemma shows. Lemma 2.14 Let the sequence (xm ) be defined by (2.22) and define the functional ψ : X → R by ψ(x) = 12 Kx − y2Y , x ∈ X. Then ψ is Fr´echet differentiable in every z ∈ X and

 ψ  (z)x = Re(Kz − y, Kx)Y = Re K ∗ (Kz − y), x X , x ∈ X . (2.23) The linear functional ψ  (z) can be identified with K ∗ (Kz −y) ∈ X in the Hilbert space X over the field R. Therefore, xm = xm−1 − a K ∗ (Kxm−1 − y) is the steepest descent step with stepsize a. Proof:

The binomial formula yields ψ(z + x) − ψ(z) − Re(Kz − y, Kx)Y =

1 Kx2Y 2

and thus   2 ψ(z + x) − ψ(z) − Re(Kz − y, Kx)Y  ≤ 1 K2 L(X,Y ) xX , 2

2.3 Landweber Iteration

43

which proves that the mapping x → Re(Kz − y, Kx)Y is the Fr´echet derivative of ψ at z.  Equation (2.22) is a linear recursion formula for xm . By induction with respect to m, it is easily seen that xm has the explicit form xm = Rm y, where the operator Rm : Y → X is defined by Rm := a

m−1 

(I − aK ∗ K)k K ∗

for m = 1, 2, . . . .

(2.24)

k=0

Choosing a singular system {μj , xj , yj : j ∈ N} for the compact and injective operator K, we see that Rm y has the representation Rm y

= a

=

=

∞ 

j=1 ∞  j=1 ∞ 

μj

m−1 

(1 − aμ2j )k (y, yj )Y xj

k=0

1 1 − (1 − aμ2j )m (y, yj )Y xj μj

q(m, μj ) (y, yj )Y xj , μj n=0

y ∈ Y,

(2.25)

with q(m, μ) = 1 − (1 − aμ2 )m . We studied this filter function q in Theorem 2.8, part (b), when we define α = 1/m. Therefore, applications of Theorems 2.6 and 2.7 yield the following result. Theorem 2.15 Again let K : X → Y be a compact and injective operator and let 0 < a < 1/K2L(X,Y ) . Let x∗ ∈ X be the exact solution of Kx∗ = y ∗ . Furthermore, let y δ ∈ Y with y δ − y ∗ Y ≤ δ. (a) Define the linear and bounded operators Rm : Y → X by (2.24). These operators Rm define a regularization strategy√with discrete regularization parameter α = 1/m, m ∈ N, and Rm L(Y,X) ≤ a m. The sequence xm,δ = Rm y δ is computed by the iteration (2.22), that is, x0,δ = 0

and

xm,δ = (I − a K ∗ K)xm−1,δ + a K ∗ y δ

(2.26)

for m = 1, 2, . . . . Every strategy m(δ) → ∞ (δ → 0) with δ 2 m(δ) → 0 (δ → 0) is admissible. (b) Again let x∗ = K ∗ z ∈ R(K ∗ ) with zY ≤ E and 0 < c1 < c2 . For every choice m(δ) with c1 Eδ ≤ m(δ) ≤ c2 Eδ , the following estimates hold: √ (2.27a) xm(δ),δ − x∗ X ≤ c3 δ E ,

 m(δ),δ ∗ Kx − y Y ≤ 1 + 1/(ac1 ) δ , (2.27b) for some c3 depending on c1 , c2 , and a. Therefore, the Landweber iteration is optimal under the information (K ∗ )−1 x∗ Y ≤ E.

44

Regularization Theory



(c) For some σ > 0, let x∗ = (K ∗ K)σ/2 z ∈ R (K ∗ K)σ/2 with zX ≤ E and let 0 < c1 < c2 . For every choice m(δ) with c1 (E/δ)2/(σ+1) ≤ m(δ) ≤ c2 (E/δ)2/(σ+1) , we have xm(δ),δ − x∗ X Kx

m(δ),δ

∗

− y Y

≤

c3 δ σ/(σ+1) E 1/(σ+1) ,

(2.27c)

≤

c3 δ ,

(2.27d)

for some c3 depending on c1 , c2 , σ, and a. Therefore, the Landweber iteration is also optimal for the information (K ∗ K)−σ/2 x∗ X ≤ E for every σ > 0. Proof: Combining the fundamental estimates (2.4a), (2.12a), and (2.12b) with Theorem 2.8 (part (b)) yields the error estimates xm,δ − x∗ X

≤

Kxm,δ − y ∗ Y

≤

√ 1 zY , am + √ 2am 1 zY , δ + am δ

for part (b) and xm,δ − x∗ X

≤

Kxm,δ − y ∗ Y

≤

 σ σ/2 am + zX , 2am  (σ+1)/2 σ+1 δ + zX , 2am δ

√

(2.28)

for part (c). Replacing m in the first term by the upper bound and in the second by the lower bound yields estimates (2.27a), (2.27b) and (2.27c), (2.27d), respectively.  The choice x0 = 0 is made to simplify the analysis. In general, the explicit iteration xm is given by xm = a

m−1 

(I − aK ∗ K)k K ∗ y + (I − aK ∗ K)m x0 ,

m = 1, 2, . . . .

k=0

In this case, Rm is affine linear, that is, of the form Rm y = z m + Sm y, y ∈ Y , for some z m ∈ X and some linear operator Sm : Y → X. For this method, we observe again that high precision (ignoring the presence of errors) requires a large number m of iterations but stability forces us to keep m small enough. We come back to the Landweber iteration in the next chapter, where we show that an optimal choice of m(δ) can be made a posteriori through a proper stopping rule. Other possibilities for regularizing first kind equations Kx = y with compact operators K, which we have not discussed, are methods using positivity or more general convexity constraints (see [27, 30, 235, 236]).

2.4 A Numerical Example

2.4

45

A Numerical Example

In this section, we demonstrate the regularization methods by Tikhonov and Landweber for the following integral equation of the first kind:

1 (1 + ts) ets x(s) ds = et ,

0 ≤ t ≤ 1,

(2.29)

0

with unique solution x∗ (t) = 1 (see Problem 2.1). The operator K : L2 (0, 1) → L2 (0, 1) is given by

1 (Kx)(t) = (1 + ts) ets x(s) ds 0 ∗

and is self-adjoint, that is, K = K. We note that x∗ does not belong to the range of K (see Problem 2.1). For the numerical evaluation of Kx, we use Simpson’s rule. With ti = i/n, i = 0, . . . , n, n even, we replace (Kx)(ti ) by ⎧ 1 n ⎨ 3n , j = 0 or n,  ti tj 4 , j = 1, 3, . . . , n − 1, wj (1 + ti tj ) e x(tj ) where wj = ⎩ 3n 2 j=0 3n , j = 2, 4, . . . , n − 2. We note that the corresponding matrix A is not symmetric. This leads  to the discretized Tikhonov equation α xα,δ + A2 xα,δ = A y δ . Here, y δ = yiδ ∈ Rn+1 is a perturbation (uniformly distributed random vector) of the discrete righthand yi∗ = exp(i/n) such that ! " n " 1  ∗ δ |y − y |2 := # (y ∗ − yiδ )2 ≤ δ . n + 1 i=0 i The average results of ten computations are given in the following tables, where we have listed the discrete norms |1 − xα,δ |2 of the errors between the exact solution x∗ (t) = 1 and Tikhonov’s approximation xα,δ (Table 2.1).

Table 2.1: Tikhonov regularization for δ = 0 α 10−1 10−2 10−3 10−4 10−5 10−6 10−7 10−8 10−9 10−10

n=8 2.4 ∗ 10−1 7.2 ∗ 10−2 2.6 ∗ 10−2 1.3 ∗ 10−2 2.6 ∗ 10−3 9.3 ∗ 10−4 3.5 ∗ 10−4 1.3 ∗ 10−3 1.6 ∗ 10−3 3.9 ∗ 10−3

n = 16 2.3 ∗ 10−1 6.8 ∗ 10−2 2.4 ∗ 10−2 1.2 ∗ 10−2 2.3 ∗ 10−3 8.7 ∗ 10−4 4.4 ∗ 10−4 3.2 ∗ 10−5 9.3 ∗ 10−5 2.1 ∗ 10−4

46

Regularization Theory Table 2.2: Tikhonov regularization for δ > 0 α 10−1 10−2 10−3 10−4 10−5 10−6 10−7 10−8 10−9 10−10

δ = 0.0001 0.2317 0.0681 0.0238 0.0119 0.0031 0.0065 0.0470 0.1018 0.1730 1.0723

δ = 0.001 0.2317 0.0677 0.0240 0.0127 0.0168 0.0909 0.2129 0.8119 1.8985 14.642

δ = 0.01 0.2310 0.0692 0.0268 0.1172 0.2553 0.6513 2.4573 5.9775 16.587

δ = 0.1 0.2255 0.1194 0.1651 1.0218 3.0065 5.9854 30.595

In the first table, we have chosen δ = 0; that is, only the discretization error for Simpson’s rule is responsible for the increase of the error for small α. This difference between discretization parameters n = 8 and n = 16 is noticeable for α ≤ 10−8 . We refer to [267] for further examples (Table 2.2). In the second table, we always took n = 16 and observed that the total error first decreases with decreasing α up to an optimal value and then increases again. This is predicted by the theory, in particular by estimates (2.21a) and (2.21b). In the following table, we list results corresponding to the iteration steps for Landweber’s method with parameter a = 0.5 and again n = 16 (Table 2.3). Table 2.3: Landweber iteration m 1 2 3 4 5 6 7 8 9 10 100 200 300 400 500

δ = 0.0001 0.8097 0.6274 0.5331 0.4312 0.3898 0.3354 0.3193 0.2905 0.2838 0.2675 0.0473 0.0248 0.0242 0.0241 0.0239

δ = 0.001 0.8097 0.6275 0.5331 0.4311 0.3898 0.3353 0.3192 0.2904 0.2838 0.2675 0.0474 0.0248 0.0242 0.0241 0.0240

δ = 0.01 0.8088 0.6278 0.5333 0.4322 0.3912 0.3360 0.3202 0.2912 0.2845 0.2677 0.0476 0.0253 0.0249 0.0246 0.0243

δ = 0.1 0.8135 0.6327 0.5331 0.4287 0.3798 0.3339 0.3248 0.2902 0.2817 0.2681 0.0534 0.0409 0.0347 0.0385 0.0424

We observe that the error decreases quickly in the first few steps and then slows down. To compare Tikhonov’s method and Landweber’s iteration, we

2.5 The Discrepancy Principle of Morozov

47

note that the error corresponding to iteration number m has to be compared with the error corresponding to α = 1/(2m) (see the estimates in the proofs of Theorems 2.15 and 2.12). Taking this into account, we observe that both methods are comparable where precision is concerned. We note, however, that the computation time of Landweber’s method is considerably higher than for Tikhonov’s method, in particular if the error δ is small. On the other hand, Landweber’s method is stable with respect to perturbations of the right-hand side and gives very good results even for large errors δ. We refer also to Section 3.5, where these regularization methods are compared with those to be discussed in the subsequent sections for Symm’s integral equation.

2.5

The Discrepancy Principle of Morozov

The following three sections are devoted to a posteriori choices of the regularization parameter α. In this section, we study a discrepancy principle based on the Tikhonov regularization method. Throughout this section, we assume again that K : X −→ Y is a compact and injective operator between Hilbert spaces X and Y with dense range R(K) ⊂ Y . Again, we study the equation Kx = y δ for given perturbations y δ ∈ Y of the exact right-hand side y ∗ = Kx∗ . The Tikhonov regularization of this equation was investigated in Section 2.2. It corresponds to the regularization operators Rα = (αI + K ∗ K)−1 K ∗

for α > 0

that approximate the unbounded inverse of K on R(K). We have seen that xα = Rα y exists and is the unique minimum of the Tikhonov functional Jα (x) := Kx − y2Y + αx2X ,

x∈X,

α > 0.

(2.30)

More facts about the dependence on α and y are proven in the following theorem. Theorem 2.16 Let y ∈ Y , α > 0, and xα be the unique solution of the equation α xα + K ∗ Kxα = K ∗ y .

(2.31)

Then xα depends continuously on y and α. The mapping α → xα X is monotonously nonincreasing and lim xα = 0 .

α→∞

The mapping α → Kxα − yY is monotonically nondecreasing and lim Kxα = y .

α→0

If y = 0, then strict monotonicity holds in both cases.

48

Regularization Theory

Proof:

We proceed in five steps.

(i) Using the definition of Jα and the optimality of xα , we conclude that

that is, xα X

α xα 2X ≤ Jα (xα ) ≤ Jα (0) = y2Y , √ ≤ yY / α. This proves that xα → 0 as α → ∞.

(ii) We choose α > 0 and β > 0 and subtract the equations for xα and xβ : α (xα − xβ ) + K ∗ K(xα − xβ ) + (α − β)xβ = 0 .

(2.32)

Multiplication by xα − xβ yields αxα − xβ 2X + K(xα − xβ )2Y = (β − α) (xβ , xα − xβ )X .

(2.33)

From this equation, we first conclude that αxα − xβ 2X ≤ |β − α| |(xβ , xα − xβ )X | ≤ |β − α| xβ X xα − xβ X , that is,

yY αxα − xβ X ≤ |β − α| xβ X ≤ |β − α| √ . β This proves continuity of the mapping α → xα . (iii) Now let β > α > 0. From (2.33) we conclude that (xβ , xα − xβ )X is (real and) positive (if zero then xα = xβ and thus xα = xβ = 0 from (2.32). This would contradict the assumption y = 0). Therefore, xβ 2X < (xβ , xα )X ≤ xβ X xα X , that is, xβ X < xα X which proves strict monotonicity of α → xα X . (iv) We multiply the normal equation for xβ by xα − xβ . This yields

 β (xβ , xα − xβ )X + Kxβ − y, K(xα − xβ ) Y = 0 . Now let α > β. From (2.33), we see that (xβ , xα − xβ )X < 0; that is, 

0 < Kxβ − y, K(xα − xβ ) Y = (Kxβ − y, Kxα − y)Y − Kxβ − y2Y . The Cauchy–Schwarz inequality yields Kxβ − yY < Kxα − yY . (v) Finally, let ε > 0. Because the range of K is dense in Y , there exists x ∈ X with Kx − y2Y ≤ ε2 /2. Choose α0 such that α0 x2X ≤ ε2 /2. Then Kxα − y2Y ≤ Jα (xα ) ≤ Jα (x) ≤ ε2 ; that is, Kxα − yY ≤ ε for all α ≤ α0 .



Now we consider the determination of α(δ) from the discrepancy principle; see [194–196]. We compute α = α(δ) > 0 such that the corresponding Tikhonov solution xα,δ , that is, the solution of the equation α xα,δ + K ∗ Kxα,δ = K ∗ y δ ,

2.5 The Discrepancy Principle of Morozov

49

that is, the minimum of Jα,δ (x) := Kx − y δ 2Y + αx2X , satisfies the equation Kxα,δ − y δ Y = δ .

(2.34)

Note that this choice of α by the discrepancy principle guarantees that, on the one side, the error of the defect is δ and, on the other side, α is not too small. Equation (2.34) is uniquely solvable, provided δ < y δ Y because by the previous theorem lim Kxα,δ − y δ Y = y δ Y > δ

α→∞

and lim Kxα,δ − y δ Y = 0 < δ .

α→0

Furthermore, α → Kxα,δ − y δ Y is continuous and strictly increasing. Theorem 2.17 Let K : X → Y be linear, compact, and one-to-one with dense range in Y . Let Kx∗ = y ∗ with x∗ ∈X, y ∗ ∈Y , and y δ ∈ Y such that y δ − y ∗ Y ≤ δ < y δ Y . Let the Tikhonov solution xα(δ) satisfy Kxα(δ),δ − y δ Y = δ for all δ ∈ (0, δ0 ). Then (a) xα(δ),δ → x∗ for δ → 0; that is, the discrepancy principle is admissible. (b) Let x∗ = K ∗ z ∈ K ∗ (Y ) with zY ≤ E. Then √ xα(δ),δ − x∗ X ≤ 2 δ E . Therefore, the discrepancy principle is an optimal regularization strategy under the information (K ∗ )−1 x∗ Y ≤ E. Proof:

xδ := xα(δ),δ minimizes the Tikhonov functional J (δ) (x) := Jα(δ),δ (x) = α(δ)x2X + Kx − y δ 2Y .

Therefore, we conclude that α(δ)xδ 2X + δ 2

= J (δ) (xδ ) ≤ J (δ) (x∗ ) = α(δ)x∗ 2X + y ∗ − y δ 2Y ≤ α(δ)x∗ 2X + δ 2 ,

and hence xδ X ≤ x∗ X for all δ > 0. This yields the following important estimate: xδ − x∗ 2X

= xδ 2X − 2 Re(xδ , x∗ )X + x∗ 2X ≤ 2 x∗ 2X − Re(xδ , x∗ )X = 2 Re(x∗ − xδ , x∗ )X .

50

Regularization Theory

First, we prove part (b): Let x∗ = K ∗ z, z ∈ Y . Then xδ − x∗ 2X

≤

2 Re(x∗ − xδ , K ∗ z)X = 2 Re(y ∗ − Kxδ , z)Y

≤ ≤

2 Re(y ∗ − y δ , z)Y + 2 Re(y δ − Kxδ , z)Y 2 δzY + 2 δzY = 4 δzY ≤ 4 δ E.

(a) Now let x∗ ∈ X and ε > 0 be arbitrary. The range R(K ∗ ) is dense in X because K is one-to-one. Therefore, there exists x ˆ = K ∗ z ∈ R(K ∗ ) such that ∗ ˆ x − x X ≤ ε/3. Then we conclude by similar arguments as above that xδ − x∗ 2X

2 Re(x∗ − xδ , x∗ − x ˆ)X + 2 Re(x∗ − xδ , K ∗ z)X ε + 2 Re(y ∗ − Kxδ , z)Y ≤ 2 x∗ − xδ X 3 ε + 4 δ zY . ≤ 2 x∗ − xδ X 3

2 This can be rewritten as x∗ − xδ X − ε/3 ≤ ε2 /9 + 4 δ zY . Now we choose δ > 0 such that the right-hand side is less than 4ε2 /9. Taking  the square root, we conclude that x∗ − xδ X ≤ ε for this δ. ≤

The condition y δ Y > δ certainly makes sense because otherwise the righthand side would be less than the error level δ, and xδ = 0 would be an acceptable approximation to x∗ . The determination of α(δ) is thus equivalent to the problem of finding the zero of the monotonic function φ(α) := Kxα,δ − y δ 2Y − δ 2 (for fixed δ > 0). It is not necessary to satisfy the equation Kxα,δ −y δ Y = δ exactly. An inclusion of the form c1 δ ≤ Kxα,δ − y δ Y ≤ c2 δ for fixed 1 ≤ c1 < c2 is sufficient to prove the assertions of the previous theorem. The computation of α(δ) can be carried out with Newton’s method. The derivative of the mapping α → xα,δ is given by the solution of the equation d α,δ x = −xα,δ , as is easily seen by differentiating (2.31) with (αI + K ∗ K) dα respect to α.

√  In the following theorem, we prove that the order of convergence O δ is best possible for the discrepancy principle. Therefore, by the results of Example 1.20, it cannot be optimal under the information (K ∗ K)−σ/2 xX ≤ E for σ > 1. Theorem 2.18 Let K be one-to-one and compact and assume that there exists 

σ > 0 such that for every x ∈ R (K ∗ K)σ/2 with y = Kx = 0, and all sequences δn → 0 and y δn ∈ Y with y − y δn Y ≤ δn and y δn Y > δn for all n, the α(δn ) is chosen by the discrepancy Tikhonov solutions xn = xα(δn ),δn (where √ principle) converge to x faster than δn to zero, that is, 1 √ xn − xX −→ 0 δn

as n → ∞ .

Then the range R(K) has to be finite-dimensional.

(2.35)

2.5 The Discrepancy Principle of Morozov

51

Proof: We show first that the choice of α(δ) by the discrepancy principle implies the boundedness of α(δ)/δ. Abbreviating xδ := xα(δ),δ , we write for δ ≤ 13 yY 1 yY 3

=

  2 1− yY ≤ yY − 2δ 3

≤ y − y δ Y + y δ Y − 2δ ≤ y δ Y − δ = y δ Y − y δ − Kxδ Y ≤ Kxδ Y 1 δ KK ∗ (y δ − Kxδ )Y ≤ K2L(X,Y ) , = α(δ) α(δ) where we applied K to (2.31). Thus we have shown that there exists c > 0 with α(δ) ≤ cδ for all sufficiently small δ. Now we assume that dim R(K) = ∞ and construct a contradiction. Let {μj , xj , yj : j ∈ N} be a singular system of K and define x :=

1 x1 μ1

y δn := y1 + δn yn

and

with

δn := μ2n .



Then y = Kx = y1 and δn → 0 as n → ∞ and x ∈ R (K ∗ K)σ/2 for every  σ > 0 and y δn − yY = δn < 1 + δn2 = y δn Y . Therefore, the assumptions for the discrepancy principle are satisfied and thus (2.35) holds. The solution of α(δn )xn + K ∗ Kxn = K ∗ y δn is given by xn =

μ1 μn δn x1 + xn . α(δn ) + μ21 α(δn ) + μ2n

We compute xn − x = −

α(δn ) μn δn

 x1 + xn 2 α(δn ) + μ2n μ1 α(δn ) + μ1

and hence √ 1 μn δn δn 1 1 n √ x − xX ≥ . = = ≥ α(δn ) + μ2n α(δn ) + δn 1 + α(δn )/δn 1+c δn This contradicts (2.35).



 We remark that the estimate α(δ) ≤ δK2L(X,Y ) / y δ Y − δ derived in the

 previous proof suggests to use δK2L(X,Y ) / y δ Y − δ as a starting value for Newton’s method to determine α(δ)! There has been an enormous effort to modify the original discrepancy principle while still retaining optimal orders of convergence. We refer to [86, 93, 102, 209, 238].

52

Regularization Theory

2.6

Landweber’s Iteration Method with Stopping Rule

It is very natural to use the following stopping criteria, which can be implemented in every iterative algorithm for the solution of Kx = y. Let r > 1 be a fixed number. Stop the algorithm at the first occurrence of m ∈ N0 with Kxm,δ − y δ Y ≤ rδ. The following theorem shows that this choice of m is possible for Landweber’s method and leads to an admissible and even optimal regularization strategy. Theorem 2.19 Let K : X → Y be linear, compact, and one-to-one with dense range. Let Kx∗ = y ∗ and y δ ∈ Y be perturbations with y ∗ − y δ Y ≤ δ and y δ Y ≥ rδ for all δ ∈ (0, δ0 ) where r > 1 is some fixed parameter (independent of δ). Let the sequence xm,δ , m = 0, 1, 2, . . . , be determined by Landweber’s method; that is, x0,δ = 0 and

 (2.36) xm+1,δ = xm,δ + a K ∗ y δ − Kxm,δ , m = 0, 1, 2, . . . , for some 0 < a < 1/K2L(X,Y ) . Then the following assertions hold: (1) limm→∞ Kxm,δ −y δ Y = 0 for every δ > 0; that is, the following stopping rule is well-defined: Let m = m(δ) ∈ N0 be the smallest integer with Kxm,δ − y δ Y ≤ rδ. (2) δ 2 m(δ) → 0 for δ → 0, that is, this choice of m(δ) is admissible. Therefore, by the assertions of Theorem 2.15, the sequence xm(δ),δ converges to x∗ as δ tends to zero. 

(3) If x∗ = K ∗ z ∈ R(K ∗ ) or x∗ = (K ∗ K)σ/2 z ∈ R (K ∗ K)σ/2 for some σ > 0 and some z, then we have the following orders of convergence: √ xm(δ),δ − x∗ X ≤ c E δ or (2.37a) xm(δ),δ − x∗ X

≤ c E 1/(σ+1) δ σ/(σ+1) ,

(2.37b)

respectively, for some c > 0 where again E = z. This means that this choice of m(δ) is optimal for all σ > 0. Proof:

In (2.25), we showed the representation Rm y =

∞  1 − (1 − aμ2j )m

μj

j=1

(y, yj )Y xj

for every y ∈ Y and thus KRm y − y2Y =

∞  j=1

 2 (1 − aμ2j )2m (y, yj )Y  .

2.6 Landweber’s Iteration Method with Stopping Rule

53

  From 1 − aμ2j  < 1, we conclude that KRm − IL(Y ) ≤ 1. Application to y δ instead of y yields Kxm,δ − y δ 2Y =

∞   2 (1 − aμ2j )2m (y δ , yj )Y  . j=1

(1) Let ε > 0 be given. Choose j1 ∈ N with ∞ 2   δ  (y , yj )Y 2 < ε . 2 j=j +1 1

Because |1 − aμ2j |2m → 0 as m → ∞ uniformly for j = 1, . . . , j1 , we can find m0 ∈ N with j1 

1 − aμ2j

 2m  δ (y , yj )Y 2

j=1

≤

max

j=1,...,j1



1 − aμ2j

j1 2  δ  2m  (y , yj )Y 2 ≤ ε 2 j=1

for m ≥ m0 .

This implies that Kxm,δ − y δ 2Y ≤ ε2 for m ≥ m0 ; that is, the method is admissible. It is sufficient to prove assertion (2) only for the case m(δ) → ∞. We set m := m(δ) for abbreviation and derive an upper bound of m. By the choice of m(δ), we have with y ∗ = Kx∗ KRm−1 y ∗ − y ∗ Y

≥ ≥

KRm−1 y δ − y δ Y − (KRm−1 − I)(y ∗ − y δ )Y r δ − KRm−1 − IL(Y ) δ ≥ (r − 1) δ , (2.38)

and hence m (r − 1)2 δ 2

≤ m

∞ 

 2 (1 − aμ2j )2m−2 (y ∗ , yj )Y 

j=1

=

∞ 

 2 m (1 − aμ2j )2m−2 μ2j (x∗ , xj )X  .

(2.39)

j=1

We show that the series converges to zero as δ → 0. (The dependence on δ is hidden in m.) First we note that mμ2 (1 − aμ2 )2m−2 ≤ 1/a for all m ≥ 1 and all μ ≥ 0 (see Problem 2.6). Now we again split the series into a finite sum and a remaining series and estimate in the “long tail”, the expression m(1 − aμ2j )2m−2 μ2j by 1/a and note that m(1 − aμ2j )2m−2 tends to zero as m → ∞ uniformly in j ∈ {1, . . . , j1 }. This proves convergence and thus part (2). For part (3) we remind the reader of the fundamental estimate (2.4a), which we need in the following form:

54

Regularization Theory xm,δ − x∗ X ≤ δ

√ a m + Rm Kx∗ − x∗ X .

(2.40)

We restrict ourselves to the case that x∗ = (K ∗ K)σ/2 z for some σ > 0. We set again E = zX and estimate m = m(δ) from above by using (2.38) and (2.28) for y ∗ instead of y δ (that is, δ = 0) and m − 1 instead of m.  (σ+1)/2 σ+1 (r − 1) δ ≤ KRm−1 y ∗ − y ∗ Y ≤ E. 2a(m − 1) Solving for m − 1 yields an estimate of the form  2/(σ+1) E m(δ) ≤ 2(m(δ) − 1) ≤ c δ for some c which depends solely on σ, r, and a. Substituting this into the first √ term of (2.40) for m = m(δ) yields δ a m(δ) ≤ ac E 1/(σ+1) δ σ/(σ+1) . Next, we estimate the second term of (2.40). In the following estimates, we use again older’s inequality with p = σ+1 that x∗ = (K ∗ K)σ/2 z and also H¨ σ and q = σ +1. Rm Kx∗ − x∗ 2X ∞ ∞   2 = (1 − aμ2j )2m |(x∗ , xj )X |2 = (1 − aμ2j )2m μ2σ j |(z, xj )X | =

j=1 ∞ 

j=1

1/p 1/q (1 − aμ2j )2m μ2σ+2 (1 − aμ2j )2m |(z, xj )X |2 |(z, xj )X |2 j

j=1

≤

⎤1/p ⎡ ⎤1/q ⎡ ∞ ∞   ⎣ (1 − aμ2j )2m μ2σ+2 |(z, xj )X |2 ⎦ ⎣ (1 − aμ2j )2m |(z, xj )X |2 ⎦ j j=1

≤ ≤

j=1 2σ/(σ+1) y ∗ Y

KRm y ∗ − E 2/(σ+1) 2σ/(σ+1) E 2/(σ+1) KRm y δ − y δ Y + (KRm − I)(y δ − y ∗ )Y .

Now we use the stopping rule for m = m(δ) in the first term and the estimate KRm − IL(Y ) ≤ 1 in the second one. This yields Rm(δ) Kx∗ − x∗ 2X ≤ (r + 1)2σ/(σ+1) E 2/(σ+1) δ 2σ/(σ+1) . Substituting this into (2.40) for m = m(δ) yields the assertion.



It is also possible to formulate a similar stopping criterion for Morozov’s discrepancy principle. Choose an arbitrary monotonically decreasing sequence (αm ) in R with limm→∞ αm = 0. Determine m = m(δ) as the smallest integer m with Kxαm ,δ − y δ Y ≤ rδ. For details, we refer the reader to [89] or [182]. One can construct more general classes of methods through the spectral representation of the solution x∗ . Comparing the regularizer xδ of Landweber’s method with the true solution x∗ , we observe that the function φ(μ) = 1/μ, μ > 0 is approximated by the polynomial Pm (μ) = 1 − (1 − aμ2 )m /μ. It is certainly possible to choose

2.7 The Conjugate Gradient Method

55

better polynomial approximations of the function μ → 1/μ. Orthogonal polynomials are particularly useful. This leads to the ν-methods; see [21, 118], or [120]. A common feature of these methods that is very crucial in the analysis is the fact that all of the polynomials Pm are independent of y and y δ . For the important conjugate gradient algorithm discussed in the next section, this is not the case, and that makes an error analysis much more difficult to obtain.

2.7

The Conjugate Gradient Method

In this section, we study the regularizing properties of the conjugate gradient method. Because the proofs of the theorems are rather technical, we only state the results and transfer the proofs to the Appendix B. First, we recall the conjugate gradient method for least squares problems for overdetermined systems of linear equations of the form Kx = y. Here, K ∈ Rm×n and y ∈ Rm with m ≥ n are given. Because it is hopeless to satisfy all equations simultaneously, one minimizes the defect f (x) := Kx − y2 , x ∈ Rn , where  ·  denotes the Euclidean norm in Rm . Standard algorithms for solving least squares problems are the QR-algorithm or the conjugate gradient method; see [71, 106, 132]. Because we assume that the latter is known for systems of equations, we formulate it now (in Figure 2.2) for the operator equation Kx = y, where K : X → Y is a bounded, linear, and injective operator between Hilbert spaces X and Y with adjoint K ∗ : Y → X. Theorem 2.20 (Fletcher–Reeves) Let K : X → Y be a compact, linear, and injective operator between Hilbert spaces X and Y . Then the conjugate gradient method is well-defined and either stops or produces sequences (xm ), (pm ) in X with the properties

 ∇f (xm ), ∇f (xj ) X = 0 for all j = m (2.41a) and (Kpm , Kpj )Y = 0

for all j = m ;

(2.41b)

that is, the gradients are orthogonal and the directions pm are K-conjugate. Furthermore, 

∇f (xj ), K ∗ Kpm X = 0 for all j < m . (2.41c) Define again the function f (x) := Kx − y2Y = (Kx − y, Kx − y)Y ,

x∈X.

We abbreviate ∇f (x) := 2 K ∗ (Kx − y) ∈ X and note that ∇f (x) is indeed the Riesz representation (see Theorem A.23) of the Fr´echet derivative f  (x) of f at x (see Lemma 2.14). We call two elements p, q ∈ X K-conjugate if (Kp, Kq)Y = 0. If K is one-to-one, this bilinear form has the properties of an inner product on X.

56

Regularization Theory

The following theorem gives an interesting and different interpretation of the elements xm .

x0 := 0 m=0 ? K y=0? ∗

@ @

@ @

yes

- STOP

yes

- STOP

no ? p0 := −K ∗ y ?

-

tm =

(Kxm − y, Kpm )Y Kpm 2Y

xm+1 = xm − tm pm

@ @

? @ K ∗ (Kxm+1 − y) = 0 ? @ no ?

γm :=

K ∗ (Kxm+1 − y)2X K ∗ (Kxm − y)2X

pm+1 := K ∗ (Kxm+1 − y) + γm pm ? m := m + 1

Figure 2.2: The conjugate gradient method

Theorem 2.21 Let (xm ) and (pm ) be the sequences of the conjugate gradient method. Define the space Vm := span {p0 , . . . , pm }. Then we have the following equivalent characterizations of Vm : Vm

= =

  span ∇f (x0 ), . . . , ∇f (xm )   span p0 , K ∗ Kp0 , . . . , (K ∗ K)m p0

(2.42a) (2.42b)

2.7 The Conjugate Gradient Method

57

for m = 0, 1, . . . . The spaces Vm are called Krylov spaces. Furthermore, xm is the minimum of f on Vm−1 for every m ≥ 1. By this result, we can write xm in the form xm = −Pm−1 (K ∗ K)p0 = Pm−1 (K ∗ K)K ∗ y

(2.43)

with a well-defined polynomial Pm−1 ∈ Pm−1 of degree at most m − 1 (which depends itself on y). Analogously, we write the defect in the form y − Kxm

= y − KPm−1 (K ∗ K)K ∗ y = y − KK ∗ Pm−1 (KK ∗ )y = Qm (KK ∗ )y

with the polynomial Qm (t) := 1 − t Pm−1 (t) of degree m. Let {μj , xj , yj : j ∈ N} be a singular system for K. If it happens that y =

N 

αj yj ∈ WN := span {y1 , . . . , yN }

j=1

for some N ∈ N, then all iterates xm ∈ AN := span {x1 , . . . , xN } because xm = Pm−1 (K ∗ K)K ∗ y =

N 

αj Pm−1 (μ2j ) μj xj .

j=1

In this exceptional case, the algorithm terminates after at most N iterations because the dimension of AN is at most N and the gradients ∇f (xi ) are orthogonal to each other. This is the reason why the conjugate gradient method applied to matrix equations stops after finitely many iterations. For operator equations in infinite-dimensional Hilbert spaces, this method produces sequences of, in general, infinitely many elements. The following characterizations of Qm are useful. Lemma 2.22 (a) The polynomial Qm minimizes the functional H(Q) := Q(KK ∗ )y2Y and satisfies

on

{Q ∈ Pm : Q(0) = 1}

H(Qm ) = Kxm − y2Y .

(b) For k = , the following orthogonality relation holds: Qk , Q  :=

∞ 

 2 μ2j Qk (μ2j ) Q (μ2j ) (y, yj )Y  = 0 .

(2.44)

j=1

If y ∈ / span {y1 , . . . , yN } for any N ∈ N, then ·, · defines an inner product on the space P of all polynomials.

58

Regularization Theory

Without a priori information, the sequence (xm ) does not converge to the solution x of Kx = y. The images, however, do converge to y. This is the subject of the next theorem. Theorem 2.23 Let K and K ∗ be one-to-one, and assume that the conjugate gradient method does not stop after finitely many steps. Then Kxm −→ y

as

m→∞

for every y ∈ Y . We give a proof of this theorem because it is a simple conclusion of the previous lemma. Proof: Let Q ∈ Pm be an arbitrary polynomial with Q(0) = 1. Then, by the previous lemma, Kxm − y2Y = H(Qm ) ≤ H(Q) =

∞ 

 2 Q(μ2j )2 (y, yj )Y  .

(2.45)

j=1

Now let ε > 0 be arbitrary. Choose j1 ∈ N such that ∞ 2    (y, yj )Y 2 < ε 2 j=j +1 1

and choose a function R ∈ C[0, μ21 ] with R(0) = 1, R∞ ≤ 1 and R(μ2j ) = 0 for j = 1, 2, . . . , j1 . By the theorem of Weierstrass, there exist polynomials ˜ m ∞ → 0 as m → ∞. We set Q ˆm = Q ˜ m /Q ˜ m (0), which ˜ m ∈ Pm with R − Q Q ˆ is possible for sufficiently large m because R(0) = 1. Then Qm converges to R ˆ m (0) = 1. Substituting this into (2.45) yields as m → ∞ and Q 

ˆm Kxm − y2Y ≤ H Q ≤

j1        ˆ m 2 ˆ m (μ2 ) − R(μ2 )2 (y, yj )Y 2 + Q Q (y, yj )Y 2 j j ∞ ( )* + j=1 j>j1 =0

≤

ˆ m − R2 y2 + Q ∞ Y

ε2 ˆ 2 Qm ∞ . 2

This expression is less than ε2 for sufficiently large m.



Now we return to the regularization of the operator equation Kx∗ = y ∗ . The operator Pm−1 (K ∗ K)K ∗ : Y → X corresponds to the regularization operator Rα of the general theory. But this operator certainly depends on the right-hand side y. The mapping y → Pm−1 (K ∗ K)K ∗ y is therefore nonlinear. So far, we have formulated and studied the conjugate gradient method for unperturbed right-hand sides. Now we consider the situation where we know only an approximation y δ of y ∗ such that y δ − y ∗ Y ≤ δ. We apply the

2.7 The Conjugate Gradient Method

59

algorithm to y δ instead of y. This yields a sequence xm,δ and polynomials Pδm and Qδm . There is no a priori strategy m = m(δ) such that xm(δ),δ converges to x∗ as δ tends to zero; see [78]. An a posteriori choice as in the previous section, however, again leads to an optimal strategy. We stop the algorithm with the smallest m = m(δ) such that the defect Kxm,δ − y δ Y ≤ rδ, where r > 1 is some given parameter. From now on, we make the assumption that y δ is never a finite linear combination of the yj . Then, by Theorem 2.23, the defect tends to zero, and this stopping rule is well-defined. We want to show that the choice m = m(δ) leads to an optimal algorithm. The following analysis, which we learned from [121] (see also [226]), is more elementary than, for example, in [21, 181], or [182]. We carry out the complete analysis but, again, postpone the proofs to Appendix B because they are rather technical. We recall that by our stopping rule Kxm(δ),δ − y δ Y ≤ rδ < Kxm(δ)−1,δ − y δ Y .

(2.46)

The following theorem establishes the optimality of the conjugate gradient method with this stopping rule. Theorem 2.24 Assume that y ∗ = Kx∗ and y δ do not belong to the linear span of finitely many yj . Let the sequence xm(δ),δ be constructed by the conjugate ∗ gradient method with stopping  rule (2.46) for fixed parameter r > 1. Let x = ∗ σ/2 ∗ σ/2 for some σ > 0 and z ∈ X. Then there exists (K K) z ∈ R (K K) c > 0 with xm(δ),δ − x∗ X ≤ c δ σ/(σ+1) E 1/(σ+1) .

(2.47)

where E = zX . As Landweber’s method and the regularization method by the spectral cutoff (but in contrast to Tikhonov’s method), the conjugate gradient method is optimal for all σ > 0 under the a priori information (K ∗ K)−σ/2 x∗ X ≤ E. There is a much simpler implementation of the conjugate gradient method for self-adjoint positive definite operators K : X → X. For such K, there exists a unique self-adjoint positive operator A : X → X with A2 = K. Let Kx = y and set z := Ax, that is, Az = y. We apply the conjugate gradient method to the equation Ax = z without knowing z. In the process of the algorithm, only the elements A∗ z = y, Apm 2 = (Kpm , pm ), and A∗ (Axm − z) = Kxm − y have to be computed. The square root A and the quantity z do not have to be known explicitly, and the method is much simpler to implement. Actually, the conjugate gradient method presented here is only one member of a large class of conjugate gradient methods. For a detailed study of these methods in connection with ill-posed problems, we refer to [104, 124, 207, 208, 210] and, in particular, the work [122].

60

2.8

Regularization Theory

Problems

2.1 Let K : L2 (0, 1) → L2 (0, 1) be the integral operator

1 (1 + ts) ets x(s) ds,

(Kx)(t) :=

0 < t < 1.

0

Show by induction that dn (Kx)(t) = dtn

1 (n + 1 + ts) sn ets x(s) ds,

0 < t < 1, n = 0, 1, . . . .

0

Prove that K is one-to-one and that the constant functions do not belong to the range of K. 2.2 Apply Tikhonov’s method of Section 2.2 to the integral equation

t x(s) ds = y(t),

0 ≤ t ≤ 1.

0

Prove that for y ∈ H 1 (0, 1) with y(0) = 0, Tikhonov’s solution xα is given by the solution of the boundary value problem −α x (t) + x(t) = y  (t), 0 < t < 1,

x(1) = 0, x (0) = 0 .

2.3 Let K : X → Y be compact For any σ ≥ 0 let the spaces 

and one-to-one. Xσ be defined by Xσ = R (K ∗ K)σ/2 equipped with the inner product

∗ σ/2  (K K) z1 , (K ∗ K)σ/2 z2 X = (z1 , z2 )X , z1 , z2 ∈ X . σ

Prove that Xσ are Hilbert spaces and that Xσ2 is compactly embedded in Xσ1 for σ2 > σ1 . 2.4 The iterated Tikhonov regularization xm,α,δ of order m ∈ N (see [87, 153]) is iteratively defined by x0,α,δ = 0,

(α I + K ∗ K)xm+1,α,δ = K ∗ y δ + α xm,α,δ

for m = 0, 1, 2, . . . . (Note that x1,α,δ is the ordinary Tikhonov regularization.)

α m is the corresponding filter func(a) Show that q (m) (α, μ) := 1 − α+μ 2 tion. (m)

(b) Prove that this filter function leads to a regularizing operator Rα √ (m) with Rα L(Y,X) ≤ m/(2 α), and q (m) satisfies (2.9) from Theorem 2.7 with ωσ(m) (α) = c αmin{σ/2,m}

2.8 Problems

61

where c depends only on m and σ. (c) Show that the iterated Tikhonov regularization of order m is asymptotically optimal under the information (K ∗ K)−σ/2 x∗ X ≤ E for every σ ≤ 2m. 2.5 Fix y δ with y δ − y ∗ Y ≤ δ and let xα,δ be the Tikhonov solution corresponding to α > 0. The curve     f (α) Kxα,δ − y δ 2Y α → := , α > 0, g(α) xα,δ 2X in R2 is called an L-curve because it has often the shape of the letter L; see [90, 125, 127]. Show by using a singular system that f  (α) = −α g  (α). Furthermore, compute the curvature    f (α) g  (α) − g  (α) f  (α) C(α) :=

3/2 f  (α)2 + g  (α)2 and show that the curvature increases monotonically for 0 < α ≤ 1/K2L(X,Y ) . 2.6 Show that mμ2 (1 − aμ2 )2m−2 ≤

1 a

1 for all m ≥ 1 and 0 ≤ μ ≤ √ . a

Chapter 3

Regularization by Discretization In this chapter, we study a different approach to regularizing operator equations of the form Kx = y, where x and y are elements of certain function spaces. This approach is motivated by the fact that for the numerical treatment of such equations, one has to discretize the continuous problem and reduce it to a finite system of (linear or nonlinear) equations. We see in this chapter that the discretization schemes themselves are regularization strategies in the sense of Chapter 2. In Section 3.1, we study the general concept of projection methods and give a necessary and sufficient condition for convergence. Although we have in mind the treatment of integral equations of the first kind, we treat the general case where K is a linear, bounded, not necessarily compact operator between (real or complex) Banach or Hilbert spaces. Section 3.2 is devoted to Galerkin methods. As special cases, we study least squares and dual least squares methods in Subsections 3.2.1 and 3.2.2. In Subsection 3.2.3, we investigate the Bubnov– Galerkin method for the case where the operator satisfies G˚ arding’s inequality. In Section 3.3, we illustrate the Galerkin methods for Symm’s integral equation of the first kind. This equation arises in potential theory and serves as a model equation for more complicated situations. Section 3.4 is devoted to collocation methods. We restrict ourselves to the moment method in Subsection 3.4.1 and to collocation by piecewise constant functions in Subsection 3.4.2, where the analysis is carried out only for Symm’s integral equation. In Section 3.5, we present numerical results for various regularization techniques (Tikhonov, Landweber, conjugate gradient, projection, and collocation methods) tested for Dirichlet boundary value problems for the Laplacian in an ellipse. Finally, we study the Backus–Gilbert method in Section 3.6. Although not very popular among mathematicians, this method is extensively used by scientists in geophysics and other applied sciences. The general ideas of Sections 3.1 and 3.2 can also be found in, for example, [17, 168, 182, 226].

© Springer Nature Switzerland AG 2021 A. Kirsch, An Introduction to the Mathematical Theory of Inverse Problems, Applied Mathematical Sciences 120, https://doi.org/10.1007/978-3-030-63343-1 3

63

64

Regularization by Discretization

3.1

Projection Methods

First, we recall the definition of a projection operator. Definition 3.1 Let X be a normed space over the field K where K = R or K = C. Let U ⊂ X be a closed subspace. A linear bounded operator P : X → X is called a projection operator on U if • P x ∈ U for all x ∈ X and • P x = x for all x ∈ U . We now summarize some obvious properties of projection operators. Theorem 3.2 Every nontrivial projection operator satisfies P 2 = P and P L(X) ≥ 1. Proof: P 2 x = P (P x) = P x follows from P x ∈ U . Furthermore, P L(X) = 2  P L(X) ≤ P 2L(X) and P = 0. This implies P L(X) ≥ 1. In the following two examples, we introduce the most important projection operators. Example 3.3 (a) (Orthogonal projection) Let X be a pre-Hilbert space over K = R or K = C and U ⊂ X be a complete subspace. Let P x ∈ U be the best approximation to x in U ; that is, P x satisfies P x − xX ≤ u − xX

for all u ∈ U .

(3.1)

By the projection theorem (Theorem A.13 of Appendix A.1), P : X → U is linear and P x ∈ U is characterized by the abstract “normal equation” (x − P x, u)X = 0 for all u ∈ U , that is, x − P x ∈ U ⊥ . In this example, by the binomial theorem we have x2X

= P x + (x − P x)2X = P x2X + x − P x2X + 2 Re(x − P x, P x)X ≥ P x2X ,    =0

that is, P L(X) = 1. Important examples of subspaces U are spaces of splines or finite elements. (b) (Interpolation operator) Let X = C[a, b] be the space of real-valued continuous functions on [a, b] supplied with the supremum norm. Then X is a normed space over R. Let U = span{u1 , . . . , un } be an n-dimensional subspace and t1 , . . . , tn ∈ [a, b] such that the interpolation problem in U is uniquely solvable; that is, det[uj (tk )] = 0. We define P x ∈ U by the interpolant of x ∈ C[a, b] in U , i.e., u = P x ∈ U satisfies u(tj ) = x(tj ) for all j = 1, . . . , n. Then P : X → U is a projection operator.

3.1

Projection Methods

65

Examples for U are spaces of algebraic or trigonometric polynomials. As a drawback of these choices, we note that from the results of Faber (see [203]) the interpolating polynomials of continuous functions x do not, in general, converge to x as the degree of the polynomials tends to infinity. For smooth periodic functions, however, trigonometric interpolation at equidistant points converges with an optimal order of convergence. We use this fact in Subsection 3.4.2. Here, as an example, we recall the interpolation by linear splines. For simplicity, we formulate only the case where the endpoints are included in the set of interpolation points. Let a = t1 < · · · < tn = b be given points, and let U ⊂ C[a, b] be defined by U

= :=

S1 (t1 , . . . , tn )   x ∈ C[a, b] : x|[tj ,tj+1 ] ∈ P1 , j = 1, . . . , n − 1 ,

(3.2)

where P1 denotes the space of polynomials of degree at most one. Then the interpolation operator Qn : C[a, b] → S1 (t1 , . . . , tn ) is given by Qn x =

n 

x(tj ) yˆj

for x ∈ C[a, b] ,

j=1

where the basis functions yˆj ∈ S1 (t1 , . . . , tn ) are defined by ⎧ t − tj−1 ⎪ , t ∈ [tj−1 , tj ] (if j ≥ 2), ⎪ ⎪ ⎪ ⎨ tj − tj−1 tj+1 − t yˆj (t) = , t ∈ [tj , tj+1 ] (if j ≤ n − 1), ⎪ ⎪ t ⎪ j+1 − tj ⎪ ⎩ 0, t∈ / [tj−1 , tj+1 ],

(3.3)

for j = 1, . . . , n. In this example Qn L(C[a,b]) = 1 (see Problem 3.1). For general interpolation operators, Qn L(X) exceeds one and Qn L(X) does not have to be bounded with respect to n. Now we define the class of projection methods. Definition 3.4 Let X and Y be Banach spaces and K : X → Y be bounded and one-to-one. Furthermore, let Xn ⊂ X and Yn ⊂ Y be finite-dimensional subspaces of dimension n and Qn : Y → Yn be a projection operator. For given y ∈ Y , the projection method for solving the equation Kx = y is to solve the equation (3.4) Qn Kxn = Qn y for xn ∈ Xn . ˆn } Equation (3.4) reduces to a finite linear system by choosing bases {ˆ x1 , . . . , x and {ˆ y1 , . . . , yˆn } of Xn and Yn , respectively. One possibility is to represent Qn y and every Qn K x ˆj , j = 1, . . . , n, in the forms Qn y =

n  i=1

βi yˆi

and Qn K x ˆj =

n  i=1

Aij yˆi ,

j = 1, . . . , n ,

(3.5)

66

Regularization by Discretization

n with βi , Aij ∈ K. The linear combination xn = j=1 αj x ˆj solves (3.4) if and only if α = (α1 , . . . , αn ) ∈ Kn solves the finite system of linear equations n 

Aij αj = βi ,

i = 1, . . . , n ;

that is, Aα = β .

(3.6)

j=1

There is a second possibility to reduce (3.4) to a finite system which is used for Galerkin methods. Let Y ∗ be the dual space of Y with dual pairing ·, · Y ∗ ,Y . We choose elements yˆi∗ ∈ Y ∗ for i = 1, . . . , n, such that the matrix Y ∈ Kn×n given yi∗ , yˆj Y ∗ ,Y is regular. Then, representing xn again as xn =

n by Yij = ˆ ˆj equation (3.4) is equivalent to j=1 αj x n 

Aij αj = βi ,

i = 1, . . . , n ;

that is, Aα = β ,

(3.7)

and βi = ˆ yi∗ , Qn y Y ∗ ,Y .

(3.8)

j=1

where now yi∗ , Qn K x ˆj Y ∗ ,Y Aij = ˆ

The orthogonal projection and the interpolation operator from Example 3.3 lead to the following important classes of projection methods, which are studied in more detail in the next sections. Example 3.5 Let K : X → Y be bounded and one-to-one. (a) (Galerkin method) Let X and Y be pre-Hilbert spaces and Xn ⊂ X and Yn ⊂ Y be finite-dimensional subspaces with dim Xn = dim Yn = n. Let Qn : Y → Yn be the orthogonal projection. Then the projected equation Qn Kxn = Qn y is equivalent to (3.9a) (Kxn , zn )Y = (y, zn )Y for all zn ∈ Yn . x1 , . . . , x ˆn } and y1 , . . . , yˆn }. Looking for a Again let Xn = span{ˆ

n Yn = span{ˆ ˆj leads to the system solution of (3.9a) in the form xn = j=1 αj x n 

αj (K x ˆj , yˆi )Y = (y, yˆi )Y

for i = 1, . . . , n ,

(3.9b)

j=1

ˆj , yˆi )Y and βi = (y, yˆi )Y . This corresponds to or Aα = β, where Aij := (K x (3.7) with yˆi∗ = yˆi after the identification of Y ∗ with Y (Theorem A.23 of Riesz). (b) (Collocation method) Let X be a Banach space, Y = C[a, b], and K : X → C[a, b] be a bounded operator. Let a = t1 < · · · < tn = b be given points (collocation points) and Yn = S1 (t1 , . . . , tn ) be the corresponding space (3.2) of n linear splines with interpolation operator Qn y = j=1 y(tj ) yˆj . Let y ∈ C[a, b] and some n-dimensional subspace Xn ⊂ X be given. Then Qn Kxn = Qn y is equivalent to (3.10a) (Kxn )(ti ) = y(ti ) for all i = 1, . . . , n .

3.1

Projection Methods

67

If we denote by {ˆ x ˆn } a basis of Xn , then looking for a solution of (3.10a) 1, . . . , x n ˆj leads to the finite linear system in the form xn = j=1 αj x n 

αj (K x ˆj )(ti ) = y(ti ),

i = 1, . . . , n ,

(3.10b)

j=1

ˆj )(ti ) and βi = y(ti ). or Aα = β, where Aij = (K x We are particularly interested in the study of integral equations of the form b k(t, s) x(s) ds = y(t) ,

t ∈ [a, b] ,

(3.11)

a

in L2 (a, b) or C[a, b] for some continuous or weakly singular function k. (3.9b) and (3.10b) now take the form Aα = β , where x =

n

j=1

(3.12)

αj x ˆj and b b Aij

=

k(t, s) x ˆj (s) yˆi (t) ds dt a

(3.13a)

a

b βi

=

y(t) yˆi (t) dt

(3.13b)

a

for the Galerkin method, and b Aij

=

k(ti , s) x ˆj (s) ds

(3.13c)

a

βi

= y(ti )

(3.13d)

for the collocation method. Comparing the systems of equations in (3.12), we observe that the computation of the matrix elements (3.13c) is less expensive than for those of (3.13a) due to the double integration for every matrix element in (3.13a). For this reason, collocation methods are generally easier to implement than Galerkin methods. On the other hand, Galerkin methods have convergence properties of high order in weak norms (superconvergence) which are of practical importance in many cases, such as boundary element methods for the solution of boundary value problems. For the remaining part of this section, we make the following assumption.

68

Regularization by Discretization

Assumption 3.6 Let K : X → Y be a linear, bounded, and injective operator between Banach spaces, Xn ⊂ X and Yn ⊂ Y be finite-dimensional subspaces of dimension n, and Qn : Y → Yn be a projection operator. We assume that  n∈N Xn is dense in X and that Qn K|Xn : Xn → Yn is one-to-one and thus invertible. Let x ∈ X be the solution of Kx = y .

(3.14)

By xn ∈ Xn , we denote the unique solutions of the equations Qn Kxn = Qn y

(3.15)

for n ∈ N. We can represent the solutions xn ∈ Xn of (3.15) in the form xn = Rn y, where Rn : Y → Xn ⊂ X is defined by Rn :=



Qn K|Xn

−1

Qn : Y −→ Xn ⊂ X .

(3.16)

The projection method is called convergent if the approximate solutions xn ∈ Xn of (3.15) converge to the exact solution x ∈ X of (3.14) for every y ∈ R(K), that is, if −1  Qn Kx −→ x, n → ∞ , (3.17) Rn Kx = Qn K|Xn for every x ∈ X. We observe that this definition of convergence coincides with Definition 2.1 of a regularization strategy for the equation Kx = y with regularization parameter α = 1/n. Therefore, the projection method converges if and only if Rn is a regularization strategy for the equation Kx = y.  Obviously, we can only expect convergence if we require that n∈N Xn is dense in X and Qn y → y for all y ∈ R(K). But, in general, this is not sufficient for convergence if K is compact. We have to assume the following boundedness condition. Theorem 3.7 Let Assumption 3.6 be satisfied. The solution xn = Rn y ∈ Xn of (3.15) converges to x for every y = Kx if and only if there exists c > 0 such that (3.18) Rn KL(X) ≤ c for all n ∈ N . If (3.18) is satisfied the following error estimate holds: xn − xX ≤ (1 + c) min zn − xX zn ∈Xn

(3.19)

with the same constant c as in (3.18). Proof: Let the projection method be convergent. Then Rn Kx → x for every x ∈ X. The assertion follows directly from the principle of uniform boundedness (Theorem A.28 of Appendix A.3).

3.1

Projection Methods

69

Now let Rn KL(X) be bounded. The operator Rn K is a projection operator −1  onto Xn because for zn ∈ Xn , we have Rn Kzn = Qn K|Xn Qn Kzn = zn . Thus we conclude that xn − x = (Rn K − I)x = (Rn K − I)(x − zn )

for all zn ∈ Xn .

This yields xn − xX ≤ (c + 1) x − zn X

for all zn ∈ Xn  and proves (3.19). Convergence xn → x follows because n∈N Xn is dense in X.  We show now a perturbation result: It is sufficient to study the question of convergence for the “principal part” of an operator K. In particular, if the projection method converges for an operator K, then convergence and the error estimates hold also for K + C, where C is compact relative to K (that is, K −1 C is compact). Theorem 3.8 Let Assumption 3.6 hold. Let C : X → Y be a linear operator with R(C) ⊂ R(K) such that K + C is one-to-one and K −1 C is compact in X. Assume, furthermore, that the projection method converges for K, that is, that Rn Kx → x for every x ∈ X, where again −1  Qn . Rn = Qn K|Xn ˜ n (K + C)x → x for every x ∈ X, Then it converges also for K + C; that is, R where now   ˜ n = Qn (K + C)|X −1 Qn . R n

∗

∗

˜ n y ∗ ∈ Xn be the ˜n = R Let x ∈ X be the solution of (K + C)x = y ∗ and x xn = Qn y ∗ . Then solution of the corresponding projected equation Qn (K + C)˜ there exists c > 0 with   ˜ xn − x∗ X ≤ c Rn y ∗ − K −1 y ∗ X + Rn Cx∗ − K −1 Cx∗ X (3.20) for all sufficiently large n. Proof: First we note that I + K −1 C = K −1 (K + C) is one-to-one and thus, because of the compactness of K −1 C, an isomorphism from X onto itself (see xn = Qn y ∗ as (I +Rn C)˜ xn = Theorem A.36). We write the equation Qn (K +C)˜ ∗ Rn y and have (I+Rn C)(˜ xn −x∗ ) = Rn y ∗ −Rn Cx∗ −x∗ = [Rn y ∗ −K −1 y ∗ ]−[Rn Cx∗ −K −1 Cx∗ ] where we haveused that x∗ + K −1 Cx∗ = K −1 y ∗ . Now we note that the operators Rn C = Rn K K −1 C converge to K −1 C in the operator norm because Rn Kx → x for every x ∈ X and K −1 C is compact in X (see Theorem A.34, part

70

Regularization by Discretization

(d), of Appendix A.3). Therefore, by the general Theorem A.37 of Appendix A.3, we conclude that (I + Rn C)−1 exist for sufficiently large n and their operator norms are uniformly bounded by some c > 0. This yields the estimate (3.20).  The first term on the right-hand side of (3.20) is just the error of the projection method for the equation Kx = y ∗ , the second for the equation Kx = Cx∗ . By the previous Theorem 3.7, these terms are estimated by Rn y ∗ − K −1 y ∗ X ∗

Rn Cx − K

−1

∗

Cx X

≤ (1 + c) min K −1 y ∗ − zn X , zn ∈Xn

≤ (1 + c) min K −1 Cx∗ − zn X zn ∈Xn

where c is a bound of Rn KL(X) . So far, we have considered the case where the right-hand side y ∗ = Kx∗ is known exactly. Now we study the case where the right-hand side is known only approximately. We understand the operator Rn from (3.16) as a regularization operator in the sense of the previous chapter. We have to distinguish between two kinds of errors on the right-hand side. The first kind corresponds to the kind of perturbation discussed in Chapter 2. Instead of the exact right-hand side y ∗ , only y δ ∈ Y is given with y δ − y ∗ Y ≤ δ. We call this the continuous perturbation of the right-hand side. A simple application of the triangle inequality yields the following result. Theorem 3.9 Let Assumption 3.6 be satisfied and let again Rn = (Qn K|Xn )−1 Qn : Y → Xn ⊂ X as in (3.16). Let x∗ ∈ X the solution of the unperturbed equation Kx∗ = y ∗ . Furthermore, we assume that the projection method converges; that is by Theorem 3.7, Rn KL(X) are uniformly bounded with respect to n. Furthermore, let y δ ∈ Y with y δ − y ∗ Y ≤ δ and xδn = Rn y δ the solution of the projected equation Qn Kxδn = Qn y δ . Then xδn − x∗ X

≤ xδn − Rn y ∗ X + Rn y ∗ − x∗ X

≤ Rn L(Y,X) y δ − y ∗ Y + Rn Kx∗ − x∗ X . (3.21)

This estimate corresponds to the fundamental estimate from (2.4a). The first term reflects the ill-posedness of the equation: The (continuous) error δ of the right-hand side is multiplied by the norm of Rn . The second term describes the discretization error for exact data. In practice, one solves the discrete systems (3.6) or (3.7) where the coefficients βi are replaced by perturbed coefficients βiδ ∈ K; that is, one solves n  j=1

Aij αjδ = βiδ ,

i = 1, . . . , n ;

that is, Aαδ = β δ ,

(3.22)

3.2

Galerkin Methods

71

instead of Aα = β where now |β δ − β|2 =

n 

|βiδ − βi |2 ≤ δ 2 .

i=1

Recall, that the elements Aij of the matrix A ∈ Kn×n and the exact coefficients βi of β ∈ Kn are given by (3.5) or (3.8). We call this the discrete perturbation of the right-hand side. Then xδn ∈ Xn is defined by xδn =

n 

αjδ x ˆj .

j=1

In this case, the choices of basis functions x ˆj ∈ Xn and yˆj ∈ Yn (and the dual basis functions yˆi∗ ) are essential. We will also see that the condition number of A reflects the ill-conditioning of the equation Kx = y. We do not carry out the analysis for these two forms of discrete perturbations in the general case but do it only for Galerkin methods in the next section.

3.2

Galerkin Methods

In this section, we assume that X and Y are (real or complex) Hilbert spaces; K : X → Y is linear, bounded, and one-to-one; Xn ⊂ X and Yn ⊂ Y are finite-dimensional subspaces with dim Xn = dim Yn = n; and Qn : Y → Yn is the orthogonal projection operator onto Yn . Then equation Qn Kxn = Qn y reduces to the Galerkin equations (see Example 3.5) (Kxn , zn )Y = (y, zn )Y

for all zn ∈ Yn .

(3.23)

If we choose bases {ˆ x1 , . . . , x ˆn } and {ˆ y1 , . . . , yˆn } of Xn and Yn , respectively, then n this leads to a finite system for the coefficients of xn = j=1 αj x ˆj (compare with (3.9b)): n  Aij αj = βi , i = 1, . . . , n , (3.24) i=1

where ˆj , yˆi )Y Aij = (K x

and βi = (y, yˆi )Y .

(3.25)

The Galerkin method is also known as the Petrov–Galerkin method (see [216]) because Petrov was the first to consider the general situation of (3.23). The special case X = Y and Xn = Yn was studied by Bubnov in 1913 and later by Galerkin in 1915 (see [99]). For this reason, this special case is also known as the Bubnov–Galerkin method. In the case when the operator K is self-adjoint and positive definite, we will see that the Bubnov–Galerkin method coincides with the Rayleigh–Ritz method; see [221] and [228].

72

Regularization by Discretization

Theorem 3.10 Let Assumption 3.6 be satisfied and let again Rn = (Qn K|Xn )−1 Qn : Y → Xn ⊂ X as in (3.16). Let x∗ ∈ X the solution of the unperturbed equation Kx∗ = y ∗ . Furthermore, we assume that the Galerkin method converges; that is by Theorem 3.7, Rn KL(X) are uniformly bounded with respect to n. (a) Let y δ ∈ Y with y δ − y ∗ Y ≤ δ and xδn = Rn y δ the solution of the projected equation Qn Kxδn = Qn y δ . Then xδn − x∗ X ≤ Rn L(Y,X) δ + Rn Kx∗ − x∗ X . (b) Let Qn y ∗ =

n



i=1

βi yˆi and βiδ ∈ K with |β δ −β| =

let αδ ∈ Kn be the solution of Aαδ = β δ . Then, xδn − x∗ X

≤

xδn − x∗ X

≤

n

(3.26)

|βiδ − βi |2 ≤ δ and

n = j=1 αjδ x ˆj ,

i=1 with xδn

an δ + Rn Kx∗ − x∗ X , (3.27a) λn Rn L(Y,X) bn δ + Rn Kx∗ − x∗ X , (3.27b)

where an

=

   n   max  ρ x ˆ j j 

X

j=1

bn

=

:

⎧ n ⎨  max  |ρi |2 ⎩ i=1

n 

 |ρj |2 = 1 ,

(3.28a)

j=1

⎫   ⎬   n  : ρ y ˆ = 1 , i i  ⎭ Y

(3.28b)

i=1

and λn > 0 denotes the smallest singular value of the matrix A. We note yi : i = that if X or Y are Hilbert spaces and {ˆ xj : j = 1, . . . , n} or {ˆ 1, . . . , n}, respectively, are orthonormal systems then an = 1 or bn = 1, respectively. Proof:

(a) has been shown before.

∗ ∗ ∗ (b) Again we use xδn − x∗ X ≤ xδn − n y X + Rn y − x X and estimate

R n ∗ ˆj with Aα = β and thus the first term. We note that Rn y = j=1 αj x

xδn − Rn y ∗ X

  n   δ  =  (α − α )ˆ x j j j  j=1

X

   n ≤ an  (αδ − α)2 j=1

= an |A−1 (β δ − β)| ≤ an |A−1 |2 |β δ − β| ≤

an δ, λn

where |A−1 |2 denotes the spectral norm of A−1 , that is, the inverse of the smallest singular value of A. This yields (3.27a).

3.2

Galerkin Methods

73

To prove (3.27b), we choose ynδ ∈ Yn such that (ynδ , yˆi )Y = βiδ for i = 1, . . . , n. Then Rn ynδ = xδn and thus xδn − Rn y ∗ X

≤

Rn L(Y,X) ynδ − Qn y ∗ Y

|(ynδ − Qn y ∗ , zn )Y | zn Y zn ∈Yn  n  δ ∗  ˆj )Y  j=1 ρj (yn − Qn y , y

n = Rn L(Y,X) sup  j=1 ρj yˆj Y ρj   n δ   j=1 ρj (βj − βj )

n = Rn L(Y,X) sup  j=1 ρj yˆj Y ρj  n 2 j=1 |ρj | δ ≤ Rn L(Y,X) |β − β| sup n ˆj Y ρj  j=1 ρj y = Rn L(Y,X) sup

≤ This ends the proof.

Rn L(Y,X) bn δ .





In the following three subsections, we derive error estimates for three special choices for the finite-dimensional subspaces Xn and Yn . The cases where Xn and Yn are coupled by Yn = K(Xn ) or Xn = K ∗ (Yn ) lead to the least squares method or the dual least squares method, respectively. Here, K ∗ : Y → X denotes the adjoint of K. In Subsection 3.2.3, we study the Bubnov–Galerkin method for the case where K satisfies G˚ arding’s inequality. In all of the subsections, we formulate the Galerkin equations for the perturbed cases first without using particular bases and then with respect to given bases in Xn and Yn .

3.2.1

The Least Squares Method

An obvious method to solve an equation of the kind Kx = y is the following: Given a finite-dimensional subspace Xn ⊂ X, determine xn ∈ Xn such that Kxn − yY ≤ Kzn − yY

for all zn ∈ Xn .

(3.29)

Existence and uniqueness of xn ∈ Xn follow easily because Xn is finite-dimensional and K is one-to-one. The solution xn ∈ Xn of this least squares problem is characterized by (Kxn , Kzn )Y = (y, Kzn )Y

for all zn ∈ Xn .

(3.30a)

We observe that this method is a special case of the Galerkin method when we set Yn := K(Xn ). Choosing a basis {ˆ xj : j = 1, . . . , n} of Xn leads to the finite system n  j=1

αj (K x ˆj , K x ˆi )Y = βi = (y, K x ˆi )Y

for i = 1, . . . , n ,

(3.30b)

74

Regularization by Discretization

or Aα = β. This has the form of (3.25) for yˆi = K x ˆi . The corresponding matrix ˆj , K x ˆi )Y is symmetric (if K = R) or Hermitian (if A ∈ Kn×n with Aij = (K x K = C) and positive definite because K is also one-to-one. Again, we study the case where the exact right-hand side y ∗ is perturbed by an error. For continuous perturbations, let xδn ∈ Xn be the solution of (Kxδn , Kzn )Y = (y δ , Kzn )Y

for all zn ∈ Xn ,

(3.31a)

where y δ ∈ Y is the perturbed right-hand side with y δ − y ∗ Y ≤ δ. ˆi )Y , i = 1, . . . , n, For the discrete perturbation, we assume that βi = (y ∗ , K x is replaced by a vector β δ ∈ Kn with |β δ − β| ≤ δ, where | · | denotes the Euclidean norm in Kn . This

nleads to the following finite system of equations ˆj : for the coefficients of xδn = j=1 αjδ x n 

αjδ (K x ˆj , K x ˆi )Y = βiδ

for i = 1, . . . , n .

(3.31b)

j=1

This system is uniquely solvable because the matrix A is positive definite. For least squares methods, the boundedness condition (3.18) is not satisfied without additional assumptions. We refer to [246] or [168], Problem 17.2, for an example. However, we can prove the following theorem. Theorem 3.11 Let K : X → Y be a linear, bounded, and injective operator between Hilbert spaces and Xn ⊂ X be finite-dimensional subspaces such that  X is dense in X. Let x∗ ∈ X be the solution of Kx∗ = y ∗ and xδn ∈ Xn n n∈N be the least squares solution from (3.31a) or (3.31b). Define   (3.32) σn := max zn X : zn ∈ Xn , Kzn Y = 1 and let there exist c > 0, independent of n, such that   min x − zn X + σn K(x − zn )Y ≤ cxX zn ∈Xn

for all x ∈ X .

(3.33)

Then the least squares method is convergent and Rn L(Y,X) ≤ σn . In this case, we have the error estimate   (3.34) x∗ − xδn X ≤ bn σn δ + c˜ min x∗ − zn X : zn ∈ Xn for some c˜ > 0. Here, bn = 1 if xδn ∈ Xn solves (3.31a); that is, δ measures the ∗ Y . If δ measures the discrete error |β δ − β| in continuous perturbation y δ − y n δ ˆj ∈ Xn , where αδ solves (3.31b), then the Euclidean norm and xn = j=1 αjδ x bn is given by ⎧ ⎫  ! n n ⎨ ⎬      K |ρj |2 :  ρ x ˆ = 1 . (3.35) bn = max  j j   ⎩ ⎭ Y j=1

j=1

If {ˆ xj : j = 1, . . . , n} is an orthonormal basis of Xn then bn = σn .

3.2

Galerkin Methods

75

Proof: It is the aim to apply Theorem 3.10. First we prove that Rn KL(X) are bounded uniformly in n. Let x ∈ X and xn := Rn Kx. Then xn satisfies (Kxn , Kzn )Y = (Kx, Kzn )Y for all zn ∈ Xn . This yields   K(xn − zn )2Y = K(xn − zn ), K(xn − zn ) Y   = K(x − zn ), K(xn − zn ) Y ≤ K(x − zn )Y K(xn − zn )Y and thus K(xn − zn )Y ≤ K(x − zn )Y for all zn ∈ Xn . Using this and the definition of σn , we conclude that xn − zn X ≤ σn K(xn − zn )Y ≤ σn K(x − zn )Y , and thus xn X

≤ ≤

xn − zn X + zn − xX + xX   xX + zn − xX + σn K(x − zn )Y .

This holds for all zn ∈ Xn . Taking the minimum, we have by Assumption (3.33) that xn X ≤ (1 + c)xX . Thus the boundedness condition (3.18) is satisfied. The application of Theorem 3.7 proves convergence. Analogously we prove the estimate for Rn L(Y,X) . Let y ∈ Y and set xn := Rn y. Then from (3.30a), we have Kxn 2Y = (Kxn , Kxn )Y = (y, Kxn )Y ≤ yY Kxn Y and thus xn X ≤ σn Kxn Y ≤ σn yY . This proves the estimate Rn L(Y,X) ≤ σn . The error estimates (3.34) follow directly from Theorem 3.10 and the estiˆj .  mates (3.26) and (3.27b) for yˆj = K x For further numerical aspects of least squares methods, we refer to [79, 82, 107, 150, 188, 189, 201, 202].

3.2.2

The Dual Least Squares Method

As the next example for the Galerkin method, we study the dual least squares method. We will see that the boundedness condition (3.18) is always satisfied. We assume in addition to the general assumptions of this section that the range R(K) of K is dense in Y . Given any finite-dimensional subspace Yn ⊂ Y , determine un ∈ Yn such that (KK ∗ un , zn )Y = (y, zn )Y

for all zn ∈ Yn ,

(3.36)

where K ∗ : Y → X denotes the adjoint of K. Then xn := K ∗ un is called the dual least squares solution. It is a special case of the Galerkin method when we set Xn := K ∗ (Yn ). Writing equation (3.36) for y = Kx in the form (K ∗ un , K ∗ zn )X = (x, K ∗ zn )X

for all zn ∈ Yn ,

76

Regularization by Discretization

we observe that the dual least squares method is just the least squares method for the equation K ∗ u = x. This explains the name. We assume again that the exact right-hand side is perturbed. Let y δ ∈ Y with y δ − y ∗ Y ≤ δ. Instead of equation (3.36), one determines xδn = K ∗ uδn ∈ Xn with (3.37) (K ∗ uδn , K ∗ zn )X = (y δ , zn )Y for all zn ∈ Yn . For discrete perturbations, we choose a basis {ˆ yj : j = 1, . . . , n} of Yn and assume that the right-hand sides βi = (y ∗ , yˆi )Y , i = 1, . . . , n, of the Galerkin equations are perturbed by a vector β δ ∈ Kn with |β δ −β| ≤ δ where |·| denotes the Euclidean norm in Kn . Instead of (3.36), we determine xδn = K ∗ uδn =

n 

αjδ K ∗ yˆj ,

j=1

where αδ ∈ Kn solves n 

αjδ (K ∗ yˆj , K ∗ yˆi )X = βiδ

for i = 1, . . . , n .

(3.38)

j=1

First we show that equations (3.37) and (3.38) are uniquely solvable. K ∗ : Y → X is one-to-one because the range R(K) is dense in Y . Thus the dimensions of Yn and Xn coincide and K ∗ is an isomorphism from Yn onto Xn . It is sufficient to prove the uniqueness of a solution to (3.37). Let un ∈ Yn with (K ∗ un , K ∗ zn )X = 0 for all zn ∈ Yn . For zn = un we conclude that 0 = (K ∗ un , K ∗ un )X = K ∗ un 2X , that is, K ∗ un = 0 or un = 0. Convergence and error estimates are proven in the following theorem. Theorem 3.12 Let X and Y be Hilbert spaces and K : X → Y be linear, bounded, and one-to-one such that the range  R(K) is dense in Y . Let Yn ⊂ Y be finite-dimensional subspaces such that n∈N Yn is dense in Y . Let x∗ ∈ X be the solution of Kx∗ = y ∗ . Then the Galerkin equations (3.37) and (3.38) are uniquely solvable for every right-hand side and every n ∈ N. The dual least squares method is convergent and   Rn L(Y,X) ≤ σn := max zn Y : zn ∈ Yn , K ∗ zn X = 1 . (3.39) Furthermore, we have the error estimates   x∗ − xδn X ≤ bn σn δ + c min x∗ − zn X : zn ∈ K ∗ (Yn )

(3.40)

xδn

for some c > 0. Here, bn = 1 if ∈ Xn solves (3.37); that is, δ measures δ ∗ − y  in Y . If δ measures the discrete error |β δ − β| and xδn = the norm y Y

n δ ∗ δ ˆj ∈ Xn , where α solves (3.38), then bn is given by j=1 αj K y ⎧ ⎫   n n ⎨ ⎬      bn = max  |ρj |2 :  ρj yˆj  =1 . (3.41)   ⎩ ⎭ Y j=1

j=1

yj : j = 1, . . . , n} forms an orthonormal system in Y . We note that bn = 1 if {ˆ

3.2

Galerkin Methods

77

Proof: We have seen already that (3.37) and (3.38) are uniquely solvable for every right-hand side and every n ∈ N. Now we prove the estimate Rn KL(X) ≤ 1, that is condition (3.18) with c = 1. Let x ∈ X and set xn := Rn Kx ∈ Xn . Then xn = K ∗ un , and un ∈ Yn satisfies (K ∗ un , K ∗ zn )X = (Kx, zn )Y for all zn ∈ Yn . For zn = un this implies xn 2X = K ∗ un 2X = (Kx, un )Y

= (x, K ∗ un )X ≤ xX xn X ,

which proves the desired estimate. If we replace Kx by y in the preceding arguments, we have xn 2X ≤ yY un Y ≤ σn yY K ∗ un X = σn yY xn X , which proves (3.39).  Finally, we show that n∈N Xn is dense in X. Let x ∈ X and ε > 0. Because ∗ ∗ K  (Y ) is dense in X, there exists y ∈ Y with x − K yX < ε/2. Because n∈N Yn is dense in Y , there exists yn ∈ Yn with y − yn Y < ε/(2KL(X,Y ) ). The triangle inequality yields that for xn := K ∗ yn ∈ Xn , x − xn X ≤ x − K ∗ yX + K ∗ (y − yn )X ≤ ε . The application of Theorem 3.10 and the estimates (3.26) and (3.27b) proves (3.40). 

3.2.3

The Bubnov–Galerkin Method for Coercive Operators

In this subsection, we assume that Y = X coincides, and K : X → X is a linear and bounded operator and Xn , n ∈ N, are finite-dimensional subspaces. The Galerkin method reduces to the problem of determining xn ∈ Xn such that (Kxn , zn )X = (y, zn )X

for all zn ∈ Xn .

(3.42)

This special case is called the Bubnov–Galerkin method. Again, we consider two kinds of perturbations of the right-hand side. If y δ ∈ X with y δ − y ∗ X ≤ δ is a perturbed right-hand side, then instead of (3.42) we study the equation (Kxδn , zn )X = (y δ , zn )X

for all zn ∈ Xn .

(3.43)

The other possibility is to choose a basis {ˆ xj : j = 1, . . . , n} of Xn and assume ˆi )X , i = 1, . . . , n, of the Galerkin equations that the right-hand sides βi = (y ∗ , x are perturbed by a vector β δ ∈ Kn with |β δ − β| ≤ δ, where | · | denotes again the Euclidean norm in Kn . In this case, instead of (3.42), we have to solve n  j=1

αjδ (K x ˆj , x ˆi )X = βiδ

for i = 1, . . . , n ,

(3.44)

78

Regularization by Discretization

for αδ ∈ Kn and set xδn =

n

j=1

αjδ x ˆj .

Before we prove a convergence result for this method, we briefly describe the Rayleigh–Ritz method and show that it is a special case of the Bubnov–Galerkin method. Let K : X → X also be self-adjoint and positive definite, that is, (Kx, y)X = (x, Ky)X and (Kx, x)X > 0 for all x, y ∈ X with x = 0. We define the functional ψ(z) := (Kz, z)X − 2 Re(y, z)X

for z ∈ X .

(3.45)

From the equation ψ(z) − ψ(x) = 2 Re(Kx − y, z − x)X + (K(z − x), z − x)X

(3.46)

and the positivity of K, we easily conclude (see Problem 3.2) that x ∈ X is the unique minimum of ψ if and only if x solves Kx = y. The Rayleigh–Ritz method is to minimize ψ over the finite-dimensional subspace Xn . From (3.46), we see that if xn ∈ Xn minimizes ψ on Xn , then, for zn = xn ± εun with un ∈ Xn and ε > 0, we have that 0 ≤ ψ(zn ) − ψ(xn ) = ±ε 2 Re(Kxn − y, un )X + ε2 (Kun , un )X for all un ∈ Xn . Dividing by ε > 0 and letting ε → 0 yields that xn ∈ Xn satisfies the Galerkin equation (3.42). If, on the other hand, xn ∈ Xn solves (3.42), then from (3.46) ψ(zn ) − ψ(xn ) = (K(zn − xn ), zn − xn )X ≥ 0 for all zn ∈ Xn . Therefore, the Rayleigh–Ritz method is identical to the Bubnov–Galerkin method. Now we generalize the Rayleigh–Ritz method and study the Bubnov–Galerkin method for the important class of coercive operators in Gelfand triples V ⊂ X ⊂ V  . For the definition of a Gelfand triple and coercive operators K : V  → V , we refer to Definition A.26 of Appendix A.3. We just recall that V  = { : V → K :

∈ V ∗ } is the space of anti-linear functionals on V with corresponding bounded sesquilinear form ·, · : V  × V → K which extends the inner product in X, that is, x, v = (x, v)X for all v ∈ V and x ∈ X. Coercivity of an operator K from V  into V means that there exists a constant γ > 0 with    x, Kx  ≥ γ x2V  for all x ∈ V  . This condition implies that K is an isomorphism from V  onto V ; see the remark following Definition A.26. If V is compactly imbedded in X, that is, j : V → X is compact, then K|X is a compact operator from X into itself. Therefore, in this subsection we measure the compactness by the “smoothing” of K from V  onto V rather than by a decay of the singular values. Now we can prove the main theorem about convergence of the Bubnov– Galerkin method for coercive operators in Gelfand triples.

3.2

Galerkin Methods

79

 Theorem 3.13 Let V ⊂ X ⊂ V  be a Gelfand triple, and Xn ⊂ V be finitedimensional subspaces such that n∈N Xn is dense in X. Let K : V  → V be coercive in the sense of Definition A.26 with constant γ > 0. Let x∗ ∈ X be the solution of Kx∗ = y ∗ for some y ∗ ∈ V . Then we have the following: (a) There exist unique solutions of the Galerkin equations (3.42)–(3.44). The Bubnov–Galerkin solutions xn ∈ Xn of (3.42) converge in V  with   (3.47) x∗ − xn V  ≤ c min x∗ − zn V  : zn ∈ Xn

for some c > 0. (b) Define ρn > 0 by

  ρn = max zn X : zn ∈ Xn , zn V  = 1

(3.48)

and the orthogonal projection operator Pn from X onto Xn . The Bubnov– Galerkin solutions converge in X (rather than in the weaker norm of V  ) if there exists c > 0 with c x − Pn xV  ≤ xX for all x ∈ X . (3.49) ρn In this case, we have the estimates Rn L(X) ≤ and ∗

x −

xδn X

ρ2n γ

(3.50)

"

  ρ2 ≤ c bn n δ + min x∗ − zn X : zn ∈ Xn γ

# (3.51)

for some c > 0. Here, bn = 1 if xδn ∈ Xn solves (3.43); that is, δ measures error |β δ − β| in the the norm y δ − y ∗ X in X.

n If δ δmeasures the discrete δ δ ˆj ∈ Xn , where α solves (3.44), then bn is Euclidean norm and xn = j=1 αj x given by ⎧ ⎫   n ⎨ ⎬   n   |aj |2 :  a x ˆ = 1 : . (3.52) bn = max  j j  ⎩ ⎭ j=1

j=1

X

Again, we note that bn = 1 if {ˆ xj : j = 1, . . . , n} forms an orthonormal system in X. Proof: (a) We apply Theorem 3.7 to the equations Kx = y, x ∈ V  , and Pn Kxn = Pn y, xn ∈ Xn , where we consider K as an operator from V  into V . We observe that the projection operator Pn is also bounded from V into Xn where we consider Xn as a subspace of V . This follows from the observation that on the finite-dimensional space Xn the norms ·X and ·V are equivalent and thus Pn uV ≤ c Pn uX ≤ c uX ≤ c˜ uV

for u ∈ V .

80

Regularization by Discretization

The constants c, and thus c˜, depend on  n. Because V is dense in X and X is dense in V  , we conclude that also n∈N Xn is dense in V  . To apply Theorem 3.7, we have to show that (3.42) is uniquely solvable in Xn and that Rn K : V  → Xn ⊂ V  is uniformly bounded with respect to n. Because (3.42) is a finite-dimensional quadratic system, it is sufficient to prove uniqueness. Let xn ∈ Xn satisfy (3.42) for y = 0. Because K is coercive, we have     γxn 2V  ≤  xn , Kxn  = (xn , Kxn )X  = 0 ; thus xn = 0. Now let x ∈ V  and set xn = Rn Kx. Then xn ∈ Xn satisfies (Kxn , zn )X = (Kx, zn )X

for all zn ∈ Xn .

(3.53)

Again, we conclude that     γxn 2V  ≤  xn , Kxn  = (xn , Kxn )X      = (xn , Kx)X  =  xn , Kx  ≤ KxV xn V  and thus xn V  ≤

1 1 KxV ≤ KL(V  ,V ) xV  . γ γ

Because this holds for all x ∈ V  , we conclude that Rn KL(V  ) ≤

1 KL(V  ,V ) . γ

Then the assumptions of Theorem 3.7 are satisfied for K : V  → V . (b) In this part we wish to apply Theorem 3.10. Let x ∈ X and xn = Rn Kx. Using the estimates (3.47) and (3.49), we conclude that x − xn X

≤ ≤

x − Pn xX + Pn x − xn X x − Pn xX + ρn Pn x − xn V 

≤ ≤

x − Pn xX + ρn Pn x − xV  + ρn x − xn V  x − Pn xX + ρn Pn x − xV  + c ρn min x − zn V 

≤

x − Pn xX + (c + 1) ρn Pn x − xV 

≤

2 xX + c1 xX = (2 + c1 ) xX ,

zn ∈Xn

and thus xn X ≤ xn − xX + xX ≤ (3 + c1 ) xX . Therefore, Rn KL(X) are uniformly bounded. Next, we prove the estimate of Rn in L(X). Let y ∈ X and xn = Rn y. We estimate       γxn 2V  ≤  xn , Kxn  = (xn , Kxn )X  = (xn , y)X  ≤

yX xn X ≤ ρn yX xn V 

3.3

Application to Symm’s Integral Equation of the First Kind

and thus xn X ≤ ρn xn V  ≤

81

1 2 ρ yX , γ n

which proves the estimate (3.50). The application of Theorem 3.10 yields the estimate (3.51).  From our general perturbation theorem (Theorem 3.8), we observe that the assumption of K being coercive can be weakened. It is sufficient to assume that K is one-to-one and satisfies G˚ arding’s inequality. We formulate the result in the next theorem. Theorem 3.14 The assertions of Theorem 3.13 also hold if K : V  → V is one-to-one and satisfies G˚ arding’s inequality with some compact operator C : V →V. For further reading, we refer to [204] and the monographs [17, 168, 182].

3.3

Application to Symm’s Integral Equation of the First Kind

In this section, we apply the Galerkin methods to an integral equation of the first kind that occurs in potential theory. We study the Dirichlet problem for harmonic functions, that is, solutions of the Laplace equation satisfying a boundary condition; that is, Δu = 0

in Ω,

u=f

on ∂Ω ,

(3.54)

where Ω ⊂ R2 is some bounded, simply connected region with analytic boundary ∂Ω and f ∈ C(∂Ω) is some given function. The single layer potential  1 u(x) = − φ(y) ln |x − y| ds(y), x ∈ Ω , (3.55) π ∂Ω

solves the boundary value problem (3.54) if and only if the density φ ∈ C(∂Ω) solves Symm’s equation  1 φ(y) ln |x − y| ds(y) = f (x) for x ∈ ∂Ω ; (3.56) − π ∂Ω

see [53]. It is well-known (see [141]) that in general the corresponding integral operator is not one-to-one. One has to make assumptions on the transfinite diameter of Ω; see [274]. We give a more elementary assumption in the following theorem. Theorem 3.15 Suppose there exists z0 ∈ Ω with |x − z0 | = 1 for all x ∈ ∂Ω. Then the only solution φ ∈ C(∂Ω) of Symm’s equation (3.56) for f = 0 is φ = 0; that is, the integral operator is one-to-one.

82

Regularization by Discretization

Proof: We give a more elementary proof than in [142], but we still need a few results from potential theory. From the continuity of x → |x − z0 |, we conclude that either |x − z0 | < 1 for all x ∈ ∂Ω or |x − z0 | > 1 for all x ∈ ∂Ω. Assume first that |x − z0 | < 1 for all x ∈ ∂Ω and choose a small disk A ⊂ Ω with center z0 such that |x − z| < 1 for all x ∈ ∂Ω and z ∈ A. Let φ ∈ C(∂Ω) satisfy (3.56) for f = 0 and define u by  1 φ(y) ln |x − y| ds(y) for x ∈ R2 . u(x) = − π ∂Ω

From potential theory (see [53]), we conclude that u is continuous in R2 , harmonic in R2 \ ∂Ω, and vanishes on ∂Ω. The maximum principle for harmonic functions implies that u vanishes in Ω. We show that u also vanishes in the exterior Ωe of Ω. The main part is to prove that  φ(y) ds(y) = 0 . φˆ := ∂Ω

Without loss of generality, we can assume that φˆ ≥ 0. We study the harmonic function v defined by  1 |x − z| φˆ ds(y), x ∈ Ωe , φ(y) ln v(x) := u(x) + ln |x − z| = π π |x − y| ∂Ω

for some z ∈ A. From the choice of A, we have v(x) =

φˆ ln |x − z| ≤ 0 π

for x ∈ ∂Ω .

We study the asymptotic behavior of v(x) as |x| tends to infinity. Elementary calculations show that   |x − z| 1 = 1 + x ˆ · (y − z) + O 1/|x|2 |x − y| |x| for |x| → ∞ uniformly in y ∈ ∂Ω, z ∈ A, and x ˆ := x/|x|. In particular, v(x) tends to zero as |x| tends to infinity. The maximum principle applied to v in Ωe yields that v(x) ≤ 0 for all x ∈ Ωe . From the asymptotic formula and ln(1 + ε) = ε + O(ε2 ), we conclude furthermore that    1 v(x) = x ˆ · φ(y) (y − z) ds(y) + O 1/|x|2 π|x| ∂Ω



and thus x ˆ·

φ(y) (y − z) ds(y) ≤ 0 for all |ˆ x| = 1 . ∂Ω

3.3

Application to Symm’s Integral Equation of the First Kind

This implies that



83

 φ(y) y ds(y) = z

∂Ω

φ(y) ds(y) . ∂Ω

$ Because this holds for all z ∈ A, we conclude that ∂Ω φ(y) ds(y) = 0. Now we see from the definition of v (for any fixed z ∈ A) that u(x) = v(x) → 0

as |x| → ∞ .

The maximum principle again yields u = 0 in Ωe . Finally, the jump conditions of the normal derivative of the single layer potential operator (see [53]) yield      2 φ(x) = lim ∇u x − εν(x) − ∇u x + εν(x) · ν(x) = 0 ε→0+

for x ∈ ∂Ω, where ν(x) denotes the unit normal vector at x ∈ ∂Ω directed into the exterior of Ω. This ends the proof for the case that maxx∈∂Ω |x − z0 | < 1. The case  minx∈∂Ω |x − z0 | > 1 is settled by the same arguments. Now we assume that the boundary ∂Ω has a parametrization of the form x = γ(s),

s ∈ [0, 2π] ,

for some 2π-periodic analytic function γ : [0, 2π] → R2 that satisfies |γ(s)| ˙ >0 for all s ∈ [0, 2π]. Then Symm’s equation (3.56) takes the form 1 − π

2π ψ(s) ln |γ(t) − γ(s)| ds = g(t) := f (γ(t))

for t ∈ [0, 2π]

(3.57)

0

  for the transformed density ψ(s) := φ γ(s) |γ(s)|, ˙ s ∈ [0, 2π]. For the special case where Ω is the disk with center 0 and radius a > 0, we have γa (s) = a (cos s, sin s) and thus ! t−s 1 ln |γa (t) − γa (s)| = ln a + ln 4 sin2 . (3.58) 2 2 For general boundaries, we can split the kernel in the form ! t−s 1 1 ln 4 sin2 − ln |γ(t) − γ(s)| = − + k(t, s) , π 2π 2

t = s ,

(3.59)

for some function k that is analytic for t = s. From the mean value theorem, we conclude that 1 ˙ . lim k(t, s) = − ln |γ(t)| s→t π

84

Regularization by Discretization

This implies that k has an analytic continuation onto [0, 2π] × [0, 2π]. With this, splitting the integral equation (3.57) takes the form 1 − 2π

2π 0

t−s ψ(s) ln 4 sin 2 2

!

2π ψ(s) k(t, s) ds = g(t)

ds +

(3.60)

0

for t ∈ [0, 2π]. We want to apply the results of the previous section on Galerkin methods to this integral equation. As the Hilbert space X, we choose X = L2 (0, 2π). The operators K, K0 , and C are defined by (Kψ)(t)

(K0 ψ)(t) Cψ

1 = − π

2π ψ(s) ln |γ(t) − γ(s)| ds,

(3.61a)

0

1 = − 2π

2π 0

" ! # t−s ψ(s) ln 4 sin2 − 1 ds, 2

= Kψ − K0 ψ

(3.61b)

(3.61c)

for t ∈ [0, 2π] and ψ ∈ L2 (0, 2π). First, we observe that K, K0 , and C are well-defined and compact operators in L2 (0, 2π) because the kernels are weakly singular (see Theorem A.35 of Appendix A.3). They are also self-adjoint in L2 (0, 2π). Then Symm’s equation (3.57) takes the form Kψ = g in the Hilbert space L2 (0, 2π). As finite-dimensional subspaces Xn and Yn , we choose the spaces of truncated Fourier series, that is,   n ˆ Xn = Yn = αj ψj : αj ∈ C , (3.62) j=−n

where ψˆj (t) = eijt for t ∈ [0, 2π] and j ∈ Z. The corresponding orthogonal projection operators Pn from L2 (0, 2π) into Xn = Yn are given as Pn ψ =

n 

αj ψˆj

(3.63)

j=−n

∞ where ψ = j=−∞ αj ψˆj . To investigate the mapping properties of K, we need the following technical result (see [168], Lemma 8.21). Lemma 3.16 1 2π

2π

ins

e 0

s& ds = ln 4 sin 2 %

2

'

−1/|n|, 0,

n ∈ Z, n = 0, n = 0.

(3.64)

3.3

Application to Symm’s Integral Equation of the First Kind

85

Proof: It suffices to study the case n ∈ N0 . First let n ∈ N. Integrating the geometric sum 1 + 2

n−1 

  s eijs + eins = i 1 − eins cot , 2 j=1

yields

0 < s < 2π ,

2π  ins  s e − 1 cot ds = 2π i . 2 0

Integration of the identity %   %  s &) s& d ( ins s e − 1 ln 4 sin2 = in eins ln 4 sin2 + eins − 1 cot ds 2 2 2 yields 2π

ins

e 0

1 s& ds = − ln 4 sin 2 in %

2

2π  0

 2π s , eins − 1 cot ds = − 2 n

which proves the assertion for n ∈ N. It remains to study the case where n = 0. Define 2π I := 0

% s& ds . ln 4 sin2 2

Then we conclude that 2π 2I

= 0

2π = 0

2π = 0

and thus I = 0.

s& ds + ln 4 sin 2 %

2

2π 0

% s& ds ln 4 cos2 2

% s s& ds ln 16 sin2 cos2 2 2 

 1 ln 4 sin s ds = 2 2

4π 0

% s& ds = I ln 4 sin2 2



This lemma shows that the functions ψˆn (t) := eint ,

t ∈ [0, 2π],

n ∈ Z,

(3.65)

86

Regularization by Discretization

are eigenfunctions of K0 : K0 ψˆn K0 ψˆ0

1 ˆ ψn |n| = ψˆ0 . =

for n = 0 and

(3.66a) (3.66b)

Now can prove the mapping properties of the operators. Theorem 3.17 Suppose there exists z0 ∈ Ω with |x − z0 | = 1 for all x ∈ ∂Ω. Let the operators K and K0 be given by (3.61a) and (3.61b), respectively. s (0, 2π), we denote the Sobolev spaces of order s (see Section A.4 of By Hper Appendix A). (a) The operators K and K0 can be extended to isomorphisms from s−1 s (0, 2π) onto Hper (0, 2π) for every s ∈ R. Hper −1/2

1/2

(b) The operator K0 is coercive from Hper (0, 2π) onto Hper (0, 2π). s−1 s (c) The operator C = K − K0 is compact from Hper (0, 2π) into Hper (0, 2π) for every s ∈ R.

Proof:

Let ψ ∈ L2 (0, 2π). Then ψ has the representation   αn eint with |αn |2 < ∞ . ψ(t) = n∈Z

n∈Z

From (3.66a) and (3.66b), we have (K0 ψ)(t) = α0 +

 1 αn eint |n|

n =0

and thus for any s ∈ R: K0 ψ2Hper s

= |α0 |2 + ⎡

(ψ, K0 ψ)L2



1 + n2

n =0

s 1 |αn |2 n2 ⎤

 1 |αn |2 ⎦ |n|

=

2π ⎣|α0 |2 +

≥

 −1/2 2π 1 + n2 |αn |2

n =0

n∈Z

=

2πψ2H −1/2 .

From the elementary estimate   s s   s−1  1 + n2 1 + n2 2 s−1  = 2 1 + n2 1+n ≤ ≤ 1 , 2 2 n 2 1+n

per

n = 0 ,

s−1 we see that K0 can be extended to an isomorphism from Hper (0, 2π) onto s Hper (0, 2π) and is coercive for s = 1/2. The operator C is bounded from

3.3

Application to Symm’s Integral Equation of the First Kind

87

r s Hper (0, 2π) into Hper (0, 2π) for all r, s ∈ R by Theorem A.47 of Appendix A.4. s−1 (0, 2π) into This proves part (c) and that K = K0 + C is bounded from Hper s s−1 (0, 2π) Hper (0, 2π). It remains to show that K is also an isomorphism from Hper s onto Hper (0, 2π). From the Riesz theory (Theorem A.36), it is sufficient to s−1 (0, 2π) with Kψ = 0. From K0 ψ = −Cψ and prove injectivity. Let ψ ∈ Hper r (0, 2π) for all r ∈ R, the mapping properties of C, we conclude that K0 ψ ∈ Hper r that is, ψ ∈ Hper (0, 2π) for all r ∈ R. In particular, this implies that ψ is continuous and the transformed function φ(γ(t))=ψ(t)/|γ(t)| ˙ satisfies Symm’s equation (3.56) for f = 0. The application of Theorem 3.15 yields φ=0. 

We are now in a position to apply all of the Galerkin methods of the previous section to Symm’s equation (3.56), that is, Kψ = g in L2 (0, 2π). We have seen that the convergence results require estimates of the condition numbers of K on the finite-dimensional spaces Xn and also approximation properties of Pn ψ. In Lemma A.45 of Appendix A, we show the following estimates for any r ≥ s. r ψn Hper

≤

s Pn ψ − ψHper

≤

s c nr−s ψn Hper for all ψn ∈ Xn , (3.67a) c r r ψHper for all ψn ∈ Hper (0, 2π) , (3.67b) nr−s

and all n ∈ N. Estimate (3.67a) is sometimes called the stability property (see −1 (0, 2π), [142]). From (3.67a) and the continuity of K −1 from L2 (0, 2π) into Hper we conclude that ψn L2 ≤ c n Kψn L2

for all ψn ∈ Xn .

(3.67c)

Indeed, this follows from ψn L2

≤

−1 −1 c n ψn Hper = c n K −1 Kψn Hper

≤

−1 Kψn L2 . c n K −1 L(L2 ,Hper )

Combining these estimates with the convergence results of the previous section, we have (almost) shown the following1 . r Theorem 3.18 Let ψ ∗ ∈ Hper (0, 2π) be the unique solution of (3.57), that is,

1 (Kψ )(t) := − π ∗

2π

ψ ∗ (s) ln |γ(t) − γ(s)| ds = g ∗ (t) := f (γ(t)) ,

0

r+1 (0, 2π) for some r ≥ 0. Let g δ ∈ L2 (0, 2π) for t ∈ [0, 2π] and some g ∗ ∈ Hper with g δ − g ∗ L2 ≤ δ and let Xn be defined by (3.62). 1 Note that in this theorem, ψ ∗ denotes the exact solution in accordance with the notation of the general theory. It must not be mixed up with the complex conjugate.

88

Regularization by Discretization

(a) Let ψnδ ∈ Xn be the least squares solution, that is, the solution of (Kψnδ , Kφn )L2 = (g δ , Kφn )L2

for all φn ∈ Xn

(3.68a)

or (b) Let ψnδ = K ψ˜nδ with ψ˜nδ ∈ Xn be the dual least squares solution, that is, ψ˜δ solves n

(K ψ˜nδ , Kφn )L2 = (g δ , φn )L2

for all φn ∈ Xn

(3.68b)

or (c) Let ψnδ ∈ Xn be the Bubnov–Galerkin solution, that is, the solution of (Kψnδ , φn )L2 = (g δ , φn )L2

for all φn ∈ Xn .

Then there exists c > 0 with ψnδ

∗

− ψ L2

1 r ≤ c n δ + r ψ ∗ Hper n

(3.68c)

! (3.69)

for all n ∈ N. 1/2

Proof: We apply Theorems 3.11, 3.12, and 3.14 (the latter with V = Hper (0, 2π) −1/2 and V  = Hper (0, 2π)). For the least squares method (Theorem 3.11), we have to show Assumption (3.33). By (3.67c) we have   (3.70) σn = max φn L2 : φn ∈ Xn , Kφn L2 = 1 ≤ c n , and thus, using (3.67b) for r = 0 and s = −1,   min ψ − φn L2 + σn K(ψ − φn )L2 φn ∈Xn

≤ ψ − Pn ψL2 + σn K(ψ − Pn ψ)L2 −1 −1 ≤ ψL2 + c n KL(Hper ≤ c ψL2 . ,L2 ) ψ − Pn ψHper Furthermore,   min ψ ∗ − φn L2 : φn ∈ Xn ≤ ψ ∗ − Pn ψ ∗ L2 ≤

c r ψ ∗ Hper nr

which proves the result for the least squares method. The estimate for the dual least squares method follows immediately from (3.40) and σn ≤ c n. For the Bubnov–Galerkin method, we have to estimate ρn from (3.48).   √ ρn = max φn L2 : φn ∈ Xn , φn Hper ≤ c n −1/2 = 1

3.3

Application to Symm’s Integral Equation of the First Kind

by (3.67a) for r = 0 and s = −1/2. This ends the proof.

89



It is interesting to note that different error estimates hold for discrete perturbations of the right-hand side. Let us denote by βk the right-hand sides of (3.68a), (3.68b), or (3.68c), respectively, for φk (t) = exp(ikt). Assume that β is perturbed by a vector β δ ∈ C2n+1 with |β − β δ | ≤ δ. We have to compute bn of (3.35), (3.41), and (3.52), respectively. Because the functions ψˆk (t) = exp(ikt), k = −n, . . . , n, are orthogonal, we compute bn for (3.41) and (3.52) by ⎫ ⎧   n ⎬ ⎨   n  1 ˆ  |ρj |2 :  ρ = 1 = √ . ψ max  j j  ⎭ ⎩ 2π L2 j=−n j=−n For the least squares method, however, we have to compute ⎧ ⎫   n n ⎨ ⎬    ˆ  |ρj |2 :  ρ K = 1 b2n = max ψ j 0 j  ⎩ ⎭ 2 j=−n

=

⎧ n ⎨

L

j=−n

⎫ ! ⎬  1 2 max |ρj |2 : 2π |ρ0 |2 + |ρ | = 1 j ⎩ ⎭ j2 j=−n j =0

2

=

n ; 2π

that is, for discrete perturbations of the right-hand side, the estimate (3.69) is asymptotically the same for the dual least squares method and the Bubnov– Galerkin method, while for the least squares method it has to be replaced by ! 1 r ψnδ − ψ ∗ L2 ≤ c n2 δ + r ψ ∗ Hper . (3.71) n The error estimates (3.69) are optimal under the a priori information ψ ∗ ∈  1/(r+1) r and ψ ∗ Hper ≤ 1. This is seen by choosing n ∼ 1/δ , which gives the asymptotic estimate   δ ∗ ψ ≤ c δ r/(r+1) . n(δ) − ψ L2 r Hper (0, 2π)

This is optimal by Problem 3.4. From the preceding analysis, it is clear that the convergence property !r−s / . 1 r s r : φn ∈ X n ≤ c ψ ∗ Hper , ψ ∗ ∈ Hper (0, 2π) , min ψ ∗ − φn Hper n and the stability property r s φn Hper ≤ c nr−s φn Hper ,

φn ∈ X n ,

for r ≥ s and n ∈ N, are the essential tools in the proofs. For regions Ω with nonsmooth boundaries, finite element spaces for Xn are more suitable. They

90

Regularization by Discretization

satisfy these conditions for a certain range of values of r and s (depending on the smoothness of the solution and the order of the finite elements). We refer to [65, 68, 141–143, 268] for more details and boundary value problems for more complicated partial differential equations. We refer to Problem 3.5 and Section 3.5, where the Galerkin methods are explicitly compared for special cases of Symm’s equation. For further literature on Symm’s and related integral equations, we refer to [11, 12, 22, 83, 223, 249, 275].

3.4

Collocation Methods

We have seen that collocation methods are subsumed under the general theory of projection methods through the use of interpolation operators. This requires the space Y to be a reproducing kernel Hilbert space, that is, a Hilbert space in which all the evaluation functionals y → y(t) for y ∈ Y and t ∈ [a, b] are bounded. Instead of presenting a general theory as in [202], we avoid the explicit introduction of reproducing kernel Hilbert spaces and investigate only two special, but important, cases in detail. First, we study the minimum norm collocation method. It turns out that this is a special case of a least squares method and can be treated by the methods of the previous section. In Subsection 3.4.2, we investigate a second collocation method for the important example of Symm’s equation. We derive a complete and satisfactory error analysis for two choices of ansatz functions. First, we formulate the general collocation method again and derive an error estimate in the presence of discrete perturbations of the right-hand side. Let X be a Hilbert space over the field K, Xn ⊂ X be finite-dimensional subspaces with dimXn = n, and a ≤ t1 < · · · < tn ≤ b be the collocation points. Let K : X → C[a, b] be bounded and one-to-one. Let Kx∗ = y ∗ , and assume that the collocation equations (Kxn )(ti ) = y(ti ) ,

i = 1, . . . , n ,

(3.72)

xj : j = are uniquely solvable in Xn for every right-hand side. Choosing a basis

{ˆ n ˆj 1, . . . , n} of Xn , we rewrite this as a system Aα = β, where xn = j=1 αj x and Aij = (K x ˆj )(ti ) , βi = y(ti ) . (3.73) The following main theorem is the analogue of Theorem 3.10 for collocation methods. We restrict ourselves to the important case of discrete perturbations of the right-hand side. Continuous perturbations could also be handled but are not of particular interest because point evaluation is no longer possible when the right-hand side is perturbed in the L2 -sense. This would require stronger norms in the range space and leads to the concept of reproducing kernel Hilbert spaces (see [168]).

3.4

Collocation Methods

91 (n)

(n)

Theorem 3.19 Let Kx∗ = y ∗ and let {t1  , . . . , tn } ⊂ [a, b], n ∈ N, be a sequence of collocation points. Assume that n∈N Xn is dense in X and that

n ˆj ∈ Xn , where αδ solves the collocation method converges. Let xδn = j=1 αjδ x Aαδ = β δ . Here, β δ ∈ Kn satisfies |β − β δ | ≤ δ where βi = y ∗ (ti ) and | · | denotes again the Euclidean norm in Kn . Then the following error estimate holds: !   an xδn − x∗ X ≤ c δ + inf x∗ − zn X : zn ∈ Xn , (3.74) λn where an

    n  = max  ρj x ˆj   j=1

:

X

n 

 |ρj | = 1 2

(3.75)

j=1

and λn denotes the smallest singular value of A. Proof: Again we write xδn − x∗ X ≤ xδn − xn X + xn − x∗ X , where ∗ xn = Rn y solves the collocation equation for β instead of β δ . The second term is estimated by Theorem 3.7. We estimate the first term by xδn − xn X

≤ ≤

an |αδ − α| = an |A−1 (β δ − β)| an δ. an |A−1 |2 |β δ − β| ≤ λn



Again we remark that an = 1 if {ˆ xj : j = 1, . . . , n} forms an orthonormal system in X.

3.4.1

Minimum Norm Collocation

Again, let K : X → C[a, b] be a linear, bounded, and injective operator from the Hilbert space X into the space C[a, b] of continuous functions on [a, b]. We assume that there exists a unique solution x∗ ∈ X of Kx∗ = y ∗ . Let a ≤ t1 < · · · < tn ≤ b be the set of collocation points. Solving the equations (3.72) in X is certainly not enough to specify the solution xn uniquely. An obvious choice is to determine xn ∈ X from the set of solutions of (3.72) that has a minimal L2 -norm among all solutions. Definition 3.20 xn ∈ X is called the moment solution of (3.72) with respect to the collocation points a ≤ t1 < · · · < tn ≤ b if xn satisfies (3.72) and   xn X = min zn X : zn ∈ X satisfies (3.72) . We can interpret this moment solution as a least squares solution. Because z → (Kz)(ti ) is bounded from X into K, Theorem A.23 by Riesz yields the existence of ki ∈ X with (Kz)(ti ) = (z, ki )X for all z ∈ X and i = 1, . . . , n. If, for example, X = L2 (a, b) and K is the integral operator b k(t, s) z(s) ds ,

(Kz)(t) = a

t ∈ [a, b] , z ∈ L2 (a, b) ,

92

Regularization by Discretization

with real-valued kernel k then ki ∈ L2 (a, b) is explicitly given by ki (s) = k(ti , s). Going back to the general case, we rewrite the moment equation (3.72) in the form (xn , ki )X = y(ti ) = (x, ki )X , i = 1, . . . , n . The minimum norm solution xn of the set of equations is characterized by the projection theorem (see Theorem A.13 of Appendix A.1) and is given by the solution of (3.72) in the space Xn := span{kj : j = 1, . . . , n}. Now we define the Hilbert space Y by Y := K(X) = R(K) with inner product (y, z)Y := (K −1 y, K −1 z)X

for y, z ∈ K(X) = R(K) .

We omit the simple proof of the following lemma. Lemma 3.21 Y is a Hilbert space that is continuously embedded in C[a, b]. Furthermore, K is an isomorphism from X onto Y . Now we can rewrite (3.72) in the form.     Kxn , Kki Y = y, Kki Y ,

i = 1, . . . , n .

Comparing this equation with (3.31a), we observe that (3.72) is the Galerkin equation for the least squares method with respect to Xn . Thus we have shown that the moment solution can be interpreted as the least squares solution for the operator K : X → Y . The application of Theorem 3.11 yields the following theorem. Theorem 3.22 Let K be one-to-one and {kj : j = 1, . . . , n} be linearly independent where kj ∈ X are such that (Kz)(tj ) = (z, kj )X for all z ∈ X, j = 1, . . . , n. Then there exists one and only one moment solution xn of (3.72). xn is given by n  xn = αj kj , (3.76) j=1

where α ∈ Kn solves the linear system Aα = β with Aij = (Kkj )(ti ) = (kj , ki )X (n)

and

βi = y(ti ) .

(3.77)

(n)

Let {t1 , . . . , tn } ⊂ [a, b], n ∈ N, be a sequence of collocation points such that  n∈N Xn is dense in X where (n)

Xn := span{kj

: j = 1, . . . , n} .

Then the moment method converges; that is, the moment solution xn ∈ Xn of (3.72) converges in X to the solution x∗ ∈ X of Kx∗ = y ∗ . If xδn =

3.4

Collocation Methods

n

93

(n)

αjδ kj , where αδ ∈ Kn solves Aαδ = β δ with |β − β δ | ≤ δ, then the following error estimate holds:   an (3.78) x∗ − xδn X ≤ δ + c min x∗ − zn X : zn ∈ Xn , λn j=1

where

   n (n)   an = max  ρ k j j  

:

X

j=1

n 

 |ρj |2 = 1

(3.79)

j=1

and where λn denotes the smallest singular value of A. Proof: The definition of  · Y implies that σn = 1, where σn is given by (3.32). Assumption (3.33) for the convergence of the least squares method is obviously satisfied because   min x∗ − zn X + σn K(x∗ − zn )Y ≤ x∗ X + σn x∗ X = 2 x∗ X .

zn ∈Xn

The application of Theorem 3.11 yields the assertion.



As an example, we again consider numerical differentiation. Example 3.23 Let X = L2 (0, 1) and K be defined by t

1 x(s) ds =

(Kx)(t) = 0



k(t, s) x(s) ds,

t ∈ [0, 1] ,

0

1, s ≤ t, 0, s > t. We choose equidistant nodes, that is, tj = nj for j = 0, . . . , n. The moment method is to minimize x2L2 under the restrictions that

with k(t, s) =

tj x(s) ds = y(tj ),

j = 1, . . . , n .

(3.80)

0

The solution xn is piecewise constant because it is a linear combination of the piecewise constant functions k(tj , ·). Therefore, the finite-dimensional space Xn is given by   (3.81) Xn = zn ∈ L2 (0, 1) : zn |(tj−1 ,tj ) constant, j = 1, . . . , n . ˆj (s) = k(tj , s). As basis functions ˆj of Xn we choose x

n x Then xn = j=1 αj k(tj , ·) is the moment solution, where α solves Aα = β with βi = y(ti ) and 1 Aij =

k(ti , s) k(tj , s) ds = 0

1 min{i, j} . n

94

Regularization by Discretization

It is not difficult to see that the moment solution is just the one-sided difference quotient xn (t1 ) =

1 y(t1 ), h

xn (tj ) =

 1 y(tj ) − y(tj−1 ) , h

j = 2, . . . , n ,

for h = 1/n. We have to check the assumptions of Theorem 3.22. First, K isone-to-one and {k(tj , ·) : j = 1 . . . , n} are linearly independent. The union n∈N Xn is dense in L2 (0, 1) (see Problem 3.6). We have to estimate an from (3.79), the smallest eigenvalue λn of A, and min x − zn L2 : zn ∈ Xn . n Let ρ ∈ Rn with j=1 ρ2j = 1. Using the Cauchy–Schwarz inequality, we estimate 2 1  1  n n   n  n+1 2   . ρ k(t , s) ds ≤ k(t , s) ds = tj = j j j   2 j=1 j=1 j=1 0

0

0 Thus an ≤ (n + 1)/2. It is straightforward to check that the inverse of A is given by the tridiagonal matrix ⎤ ⎡ 2 −1 ⎥ ⎢ −1 2 −1 ⎥ ⎢ ⎥ ⎢ .. .. .. A−1 = n ⎢ ⎥. . . . ⎥ ⎢ ⎣ −1 2 −1 ⎦ −1 1 We estimate the largest eigenvalue μmax of A−1 by the maximum absolute row sum μmax ≤ 4 n. This is asymptotically sharp because we can give a lower estimate of μmax by the trace formula n μmax ≥ trace(A−1 ) =

n   −1  A = (2n − 1) n ; jj j=1

that is, we have an estimate of λn of the form 1 1 ≤ λn ≤ . 4n 2n − 1 In Problem 3.6, it is shown that   1  min x − zn L2 : zn ∈ Xn ≤ x L2 . n Thus we have proven the following theorem. Theorem 3.24 The moment method for (3.80) converges. The following error estimate holds: 3 c n+1 ∗ δ x − xn L2 ≤ δ + (x∗ ) L2 2 n

3.4

Collocation Methods

95

if x∗ ∈ H 1 (0, 1). Here, δ is the discrete error on the right-hand side, that is,

n

n δ ∗ 2 2 δ δ ˆj , where αδ ∈ Rn solves Aαδ = β δ . j=1 |βj −y (tj )| ≤ δ and xn = j=1 αj x The choice Xn = S1 (t1 , . . . , tn ) of linear splines leads to the two-sided difference quotient (see Problem 3.8). We refer to [85, 201, 202] for further reading on moment collocation.

3.4.2

Collocation of Symm’s Equation

We now study the numerical treatment of Symm’s equation (3.57), that is, 1 (Kψ )(t) := − π ∗

2π

ψ ∗ (s) ln |γ(t) − γ(s)| ds = g ∗ (t)

(3.82)

0

for 0 ≤ t ≤ 2π by collocation methods. The integral operator K from (3.82) is 1 well-defined and bounded from L2 (0, 2π) into Hper (0, 2π). We assume throughout this subsection that K is one-to-one (see Theorem 3.15). Then we have seen in Theorem 3.17 that equation (3.82) is uniquely solvable in L2 (0, 2π) for every 1 g ∈ Hper (0, 2π); that is, K is an isomorphism. We define equidistant collocation points by π for k = 0, . . . , 2n − 1 . tk := k n There are several choices for the space Xn ⊂ L2 (0, 2π) of basis functions. Before xj : j ∈ Jn } ⊂ L2 (0, 2π) be arbitrary. we study particular cases, let Xn = span{ˆ xj : j ∈ Jn } Jn ⊂ Z denotes a set of indices with 2n elements. We assume that {ˆ forms an orthonormal system in L2 (0, 2π). The collocation equations (3.72) take the form 1 − π

2π

ψn (s) ln |γ(tk ) − γ(s)| ds = g ∗ (tk ) , k = 0, . . . , 2n − 1 ,

(3.83)

0

1 with ψn ∈ Xn . Let Qn : Hper (0, 2π) → Yn be the trigonometric interpolation operator into the 2n-dimensional space ' n−1 4  imt Yn := am e : am ∈ C . (3.84) m=−n

We recall some approximation properties of the interpolation operator Qn : 1 (0, 2π) → Yn . First, it is easily checked (see also Theorem A.46 of Hper Appendix A.4) that Qn is given by Qn ψ =

2n−1  k=0

ψ(tk ) yˆk

96

Regularization by Discretization

with Lagrange interpolation basis functions yˆk (t) =

n−1 1  im(t−tk ) e , 2n m=−n

k = 0, . . . , 2n − 1 .

(3.85)

From Theorem A.46 of Appendix A.4, we have also the estimates ψ − Qn ψL2

≤

1 Qn ψHper

≤

c 1 1 ψHper for all ψ ∈ Hper (0, 2π), n 1 1 c ψHper for all ψ ∈ Hper (0, 2π) .

(3.86a) (3.86b)

Now we can reformulate the collocation equations (3.83) as Qn Kψn = Qn g ∗

with ψn ∈ Xn .

(3.87)

We use the perturbation result of Theorem 3.8 again and split K into the form K = K0 + C with 1 (K0 ψ)(t) := − 2π

2π 0

"

t−s ψ(s) ln 4 sin 2 2

!

# − 1 ds .

(3.88)

Now we specify the spaces Xn . As a first example, we choose the orthonormal functions 1 x ˆj (t) = √ eijt for j = −n, . . . , n − 1 . (3.89) 2π We prove the following convergence result. Theorem 3.25 Let x ˆj , j = −n, . . . , n − 1, be given by (3.89). The collocation method is convergent; that is, the solution ψn ∈ Xn of (3.83) converges to the solution ψ ∗ ∈ L2 (0, 2π) of (3.82) in L2 (0, 2π). Let the right-hand side of (3.83) be replaced by β δ ∈ C2n with 2n−1 

|βkδ − g ∗ (tk )|2 ≤ δ 2 .

k=0

ˆj )(tk ). Then the Let αδ ∈ C2n be the solution of Aαδ = β δ , where Akj = (K x following error estimate holds:   √ ψnδ − ψ ∗ L2 ≤ c n δ + min ψ ∗ − φn L2 : φn ∈ Xn , (3.90) where

n−1 1  δ ijt ψnδ (t) = √ αj e . 2π j=−n

r If ψ ∗ ∈ Hper (0, 2π) for some r > 0, then " # √ 1 δ ∗ ∗ r n δ + r ψ Hper ψn − ψ L2 ≤ c . n

(3.91)

3.4

Collocation Methods

97

Proof: By the perturbation Theorem 3.8, it is sufficient to prove the result for K0 instead of K. By (3.66a) and (3.66b), the operator K0 maps Xn into Yn = Xn . Therefore, the collocation equation (3.87) for K0 reduces to K0 ψn = Qn g ∗ . We want to apply Theorem 3.7 and have to estimate Rn K0 where in this case −1  1 Rn = K0 |Xn Qn . Because K0 : L2 (0, 2π) → Hper (0, 2π) is invertible, we conclude that 1 1 1 Rn gL2 = ψn L2 ≤ c1 K0 ψn Hper = c1 Qn gHper ≤ c2 gHper

1 for all g ∈ Hper (0, 2π), and thus 1 Rn KψL2 ≤ c2 KψHper ≤ c3 ψL2

for all ψ ∈ L2 (0, 2π). The application of Theorem 3.7 yields convergence. To prove the error estimate (3.90), we want to apply Theorem 3.19 and hence have to estimate the singular values of the matrix B defined by ˆj )(tk ), Bkj = (K0 x

k, j = −n, . . . , n − 1 ,

with x ˆj from (3.89). From (3.66a) and (3.66b), we observe that 1 1 ijk π e n, Bkj = √ 2π |j|

k, j = −n, . . . , n − 1 ,

where 1/|j| has to be replaced by 1 if j = 0. Because the singular values of B are the square roots of the eigenvalues of B ∗ B, we compute 

B∗B

 j

=

n−1  k=−n

B k Bkj =

n−1  π 1 1 n 1 eik(j−) n = δj , 2π | ||j| π 2 k=−n

where again 1/ 2 has to be replaced 0 by 1 for = 0. From this, we see that the singular values √ of B are given by n/(π 2 ) for = 1, . . . , n. The smallest singular value is 1/ nπ. Estimate (3.74) of Theorem 3.19 yields the assertion. (3.91) follows from Theorem A.46.  Comparing the estimate (3.91) with the corresponding error estimate (3.69) for the Galerkin methods, it seems that the estimate √ for the collocation method is better because the error δ is only multiplied by n instead of n. Let us now compare the errors of the continuous perturbation y ∗ − y δ L2 with the discrete perturbation for both methods. To do this, we have to “extend” the discrete vector β δ to a function ynδ ∈ Xn . For the collocation method, we have to use

2n−1 the interpolation operator Qn and define ynδ ∈ Xn by ynδ = j=1 βjδ yˆj , where yˆj are the Lagrange basis functions (3.85). Then ynδ (tk ) = βkδ , and we estimate ynδ − y ∗ L2 ≤ ynδ − Qn y ∗ L2 + Qn y ∗ − y ∗ L2 .

98

Regularization by Discretization

Writing ynδ (t) − Qn y ∗ (t) =

n−1 

ρj eijt ,

j=−n

a simple computation shows that n−1 

|βkδ − y ∗ (tk )|2

=

k=0

n−1 

|βkδ − Qn y ∗ (tk )|2 = n−1 

2n

  

k=0 j=−n

k=0

=

n−1   n−1 

j=−n

|ρj |2 =

2 π ρj eikj n 

n δ y − Qn y ∗ 2L2 . π n

(3.92)

Therefore, for the collocation error 0 method we have to compare the continuous √ δ with the discrete error δ n/π. This gives an extra factor of n in the first terms of (3.90) and (3.91).

n 1 δ For Galerkin methods, however, we define ynδ (t) = 2π j=−n βj exp(ijt).  δ ij·  δ Then yn , e L2 = βj . Let Pn be the orthogonal projection onto Xn . In ynδ − y ∗ L2 ≤ ynδ − Pn y ∗ L2 + Pn y ∗ − y ∗ L2 , we estimate the first term as ynδ − Pn y ∗ 2L2 =

n 1   δ  ∗ ij·  2 β − y , e L2 . 2π j=−n j

In this case, the continuous and discrete errors are of the same order. Choosing trigonometric polynomials as basis functions is particularly suitable for smooth boundary data. If ∂Ω or the right-hand side f of the boundary value problem (3.54) is not smooth, then spaces of piecewise constant functions are more appropriate. We now study the case where the basis functions x ˆj ∈ L2 (0, 2π) are defined by ' 0 n x ˆ0 (t)

=

π,

0, ' 0 n

x ˆj (t)

=

π,

0,

π if t < 2n or t > 2π − π π if 2n < t < 2π − 2n ,

if |t − tj | < if |t − tj | >

π 2n , π 2n ,

π 2n ,

(3.93a)

(3.93b)

for j = 1, . . . , 2n − 1. Then x ˆj , j = 0, . . . , 2n − 1, are also orthonormal in L2 (0, 2π). In the following lemma, we collect some approximation properties of the corresponding spaces Xn .

3.4

Collocation Methods

99

Lemma 3.26 Let Xn = span{ˆ xj : j = 0, . . . , 2n − 1}, where x ˆj are defined 2 : L (0, 2π) → X be the orthogonal projection by (3.93a) and (3.93b). Let P n n  operator. Then n∈N Xn is dense in L2 (0, 2π) and there exists c > 0 with ψ − Pn ψL2

≤

K(ψ − Pn ψ)L2

≤

c 1 1 ψHper for all ψ ∈ Hper (0, 2π), n c ψL2 for all ψ ∈ L2 (0, 2π) . n

(3.94a) (3.94b)

Proof: Estimate (3.94a) is left as an exercise. To prove estimate (3.94b), we use (implicitly) a duality argument:   K(ψ − Pn ψ), φ L2 K(ψ − Pn ψ)L2 = sup φL2 φ =0   ψ − Pn ψ, Kφ L2 = sup φL2 φ =0   ψ, (I − Pn )Kφ L2 = sup φL2 φ =0 ≤ ≤

(I − Pn )KφL2 φL2 φ =0

ψL2 sup

1 KφHper c˜ ψL2 sup n φL2 φ =0

c ψL2 . n

≤



Before we prove a convergence theorem, we compute the singular values of the matrix B defined by Bkj

1 = (K0 x ˆj )(tk ) = − 2π

"

2π

2 tk

x ˆj (s) ln 4 sin 0

−s 2

!

# − 1 ds.

(3.95)

Lemma 3.27 B is symmetric and positive definite. The singular values of B coincide with the eigenvalues and are given by 3 3  n n sin mπ 1 2n and μm = (3.96a) μ0 = m 2 π π 2nπ +j j∈Z

2n

for m = 1, . . . , 2n − 1. Furthermore, there exists c > 0 with √ 1 √ ≤ μm ≤ c n πn

for all m = 0, . . . , 2n − 1 .

(3.96b)

We observe that the condition number of B, that is, the ratio between the largest and smallest singular values, is again bounded by n.

100

Regularization by Discretization

Proof:

We write Bkj = −

1 2π

3

n π

π tj + 2n "



ln 4 sin2

π tj − 2n

with

3

1 b = − 2π

n π

π t + 2n



π t − 2n

s − tk 2

!

# − 1 ds = bj−k

( % ) s& ln 4 sin2 − 1 ds , 2

where we extended the definition of t to all ∈ Z. Therefore, B is circulant and symmetric. The eigenvectors x(m) and eigenvalues μm of B are given by x(m) =

 imk π 2n−1 n e k=0

and μm =

2n−1 

π

bk eimk n ,

k=0

respectively, for m = 0, . . . , 2n − 1, as is easily checked. We write μm in the form 2π ) ( % s& 1 − 1 ds = K0 ψm (0) ψm (s) ln 4 sin2 μm = − 2π 2 0

3

with

n imk π π n e , k ∈ Z. for |s − tk | ≤ ψm (s) = π 2n

Let ψm (t) = k∈Z ρm,k exp(ikt). Then by (3.66a) and (3.66b), we have  ρm,k . μm = ρm,0 + |k| k =0

Therefore, we have to compute the Fourier coefficients ρm,k of ψm . They are given by ρm,k

1 = 2π

2π

3 −iks

ψm (s) e 0

ds =

π

tj + 2n  2n−1 n 1  imj π e n e−iks ds . π 2π j=0

For k = 0, this reduces to 3  0 2n−1 n 1  imj π n/π e n = ρm,0 = 0 π 2n j=0

π tj − 2n

if m = 0, if m = 1, . . . , 2n − 1,

and for k = 0 to 3 ρm,k =

⎧ 3 n 2n πk ⎨ πk 2n−1  sin sin , π n 2n ei(m−k)j n = π πk 2n ⎩ π πk j=0 0,

if m − k ∈ 2nZ, if m − k ∈ / 2nZ.

3.4

Collocation Methods

101

0 Thus we have μ0 = n/π and 3 3  n 2n  | sin πk n sin πm 1 2n | 2n μm = = m 2 2 π π k π 2nπ j∈Z 2n + j k−m∈2nZ for m = 1, . . . , 2n − 1. This proves (3.96a). Because all eigenvalues are positive, the matrix B is positive definite and the eigenvalues coincide with the singular values. We set x = m/(2n) ∈ (0, 1) and separate the first two terms in the series. This yields 3 ! n sin πx 1 1 + μm = (3.97) π 2nπ x2 (x − 1)2 3 3 ∞ ∞ n sin πx  1 n sin πx  1 + + π 2nπ j=1 (x + j)2 π 2nπ j=2 (x − j)2 3 ! n 1 sin πx sin π(1 − x) ≥ + . π 2nπ x2 (1 − x)2 From the elementary estimate sin πx sin π(1 − x) + ≥ 8, 2 x (1 − x)2 we conclude that μm ≥

x ∈ (0, 1) ,

1 4 1 √ ≥ √ π πn πn

for m = 1, . . . , 2n − 1. The upper estimate of (3.96b) is proven analogously. Indeed, from (3.97) we estimate, using also | sin t/t| ≤ 1, ⎛ ⎞ 3 ∞  n 1 ⎝ sin πx 1⎠ sin π(1 − x) + + 2 μm ≤ π 2nπ x2 (1 − x)2 j2 j=1 ⎛ ⎞ 3 ∞ √ 1 2 1⎠ n 1 ⎝1 ≤ + + ≤ c n π 2n x 1−x π j=1 j 2 for some c > 0.



Now we can prove the following convergence result. Theorem 3.28 Let x ˆj , j = 0, . . . , 2n − 1, be defined by (3.93a) and (3.93b). The collocation method is convergent; that is, the solution ψn ∈ Xn of (3.83) converges to the solution ψ ∗ ∈ L2 (0, 2π) of (3.82) in L2 (0, 2π). Let the right-hand side be replaced by β δ ∈ C2n with 2n−1  j=0

|βjδ − g ∗ (tj )|2 ≤ δ 2 .

102

Regularization by Discretization

Let αδ ∈ C2n be the solution of Aαδ = β δ , where Akj = K x ˆj (tk ). Then the following error estimate holds:  √  (3.98) ψnδ − ψ ∗ L2 ≤ c n δ + min ψ ∗ − φn L2 : φn ∈ Xn ,

2n−1 1 where ψnδ = j=0 αjδ x ˆj . If ψ ∗ ∈ Hper (0, 2π), then " # √ 1 1 ψ ∗ Hper nδ + . (3.99) ψnδ − ψ ∗ L2 ≤ c n Proof: By the perturbation theorem (Theorem 3.8), it is sufficient to prove the result for K0 instead of K. Again set −1  1 Qn : Hper (0, 2π) −→ Xn ⊂ L2 (0, 2π) , Rn = Qn K0 |Xn

2n−1 1 let ψ ∈ Hper (0, 2π), and set ψn = Rn ψ = j=0 αj x ˆj . Then α ∈ C2n solves Bα = β with βk = ψ(tk ), and thus by (3.96b) ψn L2 = |α| ≤ |B

−1

:1/2 92n−1  √ 2 |2 |β| ≤ πn |ψ(tk )| k=0

where | · | again denotes the Euclidean norm in Cn . Using this estimate and (3.92) for βkδ = 0, we conclude that Rn ψL2 = ψn L2 ≤ n Qn ψL2

(3.100)

1 for all ψ ∈ Hper (0, 2π). Thus

Rn K0 ψL2 ≤ n Qn K0 ψL2 for all ψ ∈ L2 (0, 2π). Now we estimate Rn K0 ψL2 by the L2 -norm of ψ itself. Let ψ˜n = Pn ψ ∈ Xn be the orthogonal projection of ψ ∈ L2 (0, 2π) in Xn . Then Rn K0 ψ˜n = ψ˜n and ψ˜n L2 ≤ ψL2 , and thus      Rn K0 ψ − ψ˜n L2 = Rn K0 (ψ − ψ˜n )L2 ≤ n Qn K0 ψ − ψ˜n L2 ≤ n Qn K0 ψ − K0 ψL2 + n K0 ψ − K0 ψ˜n L2 + n K0 ψ˜n − Qn K0 ψ˜n L2 . Now we use the error estimates (3.86a), (3.94a), and (3.94b) of Lemma 3.26. This yields ( ) 1 1 Rn K0 ψ − ψ˜n L2 ≤ c1 K0 ψHper + ψL2 + K0 ψ˜n Hper ( ) ≤ c2 ψL2 + ψ˜n L2 ≤ c3 ψL2 , that is, Rn K0 ψL2 ≤ c4 ψL2 for all ψ ∈ L2 (0, 2π). Therefore, the assumptions of Theorem 3.7 are satisfied. The application of Theorem 3.19 yields the error estimate (3.99). 

3.5

Numerical Experiments for Symm’s Equation

103

Among the extensive literature on collocation methods for Symm’s integral equation and related equations, we mention only the work of [8, 66, 67, 140, 147, 241, 242]. Symm’s equation has also been numerically treated by quadrature methods; see [90, 169, 239, 240, 250, 251]. For more general problems, we refer to [9, 68].

3.5

Numerical Experiments for Symm’s Equation

In this section, we apply all of the previously investigated regularization strategies to Symm’s integral equation 1 (Kψ)(t) := − π

2π ψ(s) ln |γ(t) − γ(s)| ds = g(t),

0 ≤ t ≤ 2π ,

0

  where in this example γ(s) = cos s, 2 sin s , 0 ≤ s ≤ 2π, denotes the parametrization of the ellipse with semiaxes 1 and 2. First, we discuss the numerical computation of Kψ. We write Kψ in the form (see (3.60)) (Kψ)(t) = −

1 2π

2π 0

2π % t − s& ds + ψ(s) ln 4 sin2 ψ(s) k(t, s) ds , 2 0

for 0 ≤ t ≤ 2π, with the analytic function k(t, s) k(t, t)

1 |γ(t) − γ(s)|2 ln , t = s, 2π 4 sin2 t−s 2 1 ˙ , 0 ≤ t ≤ 2π . = − ln |γ(t)| π = −

We use the trapezoidal rule for periodic functions (see [168]). Let tj = j nπ , j = 0, . . . , 2n − 1. The smooth part is approximated by 2π k(t, s) ψ(s) ds ≈ 0

2n−1 π  k(t, tj ) ψ(tj ), n j=0

0 ≤ t ≤ 2π .

For the weakly singular part, we replace ψ by its trigonometric interpolation

2n−1 polynomial Qn ψ = j=0 ψ(tj ) Lj into the 2n-dimensional space ⎧ n ⎨ ⎩

j=0

aj cos(jt) +

n−1  j=1

⎫ ⎬ bj sin(jt) : aj , bj ∈ R ⎭

104

Regularization by Discretization

over R (see Section A.4 of Appendix A.4). From (A.37) and Lemma 3.16, we conclude that 1 − 2π

2π 0

% t − s& ds ψ(s) ln 4 sin2 2

≈

=

1 − 2π 2n−1 

2π  0

%  t − s& Qn ψ (s) ln 4 sin2 ds 2

ψ(tj ) Rj (t) ,

0 ≤ t ≤ 2π ,

j=0

where Rj (t)

2π % t − s& 1 ds = − Lj (s) ln 4 sin2 2π 2 0 ' 4 n−1  1 1 1 = cos n(t − tj ) + cos m(t − tj ) n 2n m m=1

for j = 0, . . . , 2n − 1. Therefore, the operator K is replaced by (Kn ψ)(t) :=

2n−1  j=0

  π k(t, tj ) , ψ(tj ) Rj (t) + n

0 ≤ t ≤ 2π .

It is well-known (see [168]) that Kn ψ converges uniformly to Kψ for every 2π-periodic continuous function ψ. Furthermore, the error Kn ψ − Kψ∞ is exponentially decreasing for analytic functions ψ. For t = tk , k = 0, . . . , 2n − 1,

2n−1 we have (Kn ψ)(tk ) = j=0 Akj ψ(tj ) with the symmetric matrix Akj := R|k−j| +

where

1 R = n

'

π k(tk , tj ), n

k, j = 0 . . . , 2n − 1 ,

n−1  1 (−1) m π + cos 2n m n m=1

4 ,

= 0, . . . , 2n − 1.

  For the numerical example, we take ψ(s) = exp 3 sin s , 0 ≤ s ≤ 2π, and   ˜ We g = Kψ or, discretized, ψ˜j = exp 3 sin tj , j = 0, . . . , 2n − 1, and g˜ = Aψ. take n = 60 and add uniformly distributed random noise on the data g˜. All the results show the average of 10 computations. The errors are measured in the 2n−1 1 2 2n discrete norm |z|22 := 2n j=0 |zj | , z ∈ C . First, we consider Tikhonov’s regularization method for δ = 0.1, δ = 0.01, δ = 0.001, and δ = 0. In Figure 3.1, we plot the errors ψ˜α,δ − ψ˜2 and   α,δ Aψ˜ − g˜ in the solution and the right-hand side, respectively, versus the 2 regularization parameter α. We clearly observe the expected behavior of the errors: For δ > 0 the error in the solution has a well-defined minimum that depends on δ, while the defect always converges to zero as α tends to zero.

3.5

Numerical Experiments for Symm’s Equation

105

The minimal values errδ of the errors in the solution are approximately 0.351, 0.0909, and 0.0206 for δ = 0.1, 0.01, and 0.001, respectively. From this, we observe the order of convergence: increasing the error by factor 10 should increase the error by factor 102/3 ≈ 4.64, which roughly agrees with the numerical results where errδ=0.1 /errδ=0.01 ≈ 3.86 and errδ=0.01 /errδ=0.001 ≈ 4.41. In Figure 3.2, we show the results for the Landweber iteration with a = 0.5 for the same example where again δ = 0.1, δ = 0.01, δ = 0.001, and δ = 0. The errors in the solution and the defects are now plotted versus the iteration number m.

1

0.5

0.9

0.45

0.8

0.4

0.7

0.35

0.6

0.3

0.5

0.25

0.4

0.2

0.3

0.15

0.2

0.1

0.1

0.05

0 0

0.01

0.02

0.03

0.04

0.05

0.06

0.07

0.04

0.08

0 0

0.005

0.01

0.015

0.02

0.025

0.03

0.035

0.04

0.02 0.018

0.035

0.016 0.03 0.014 0.025

0.012

0.02

0.01 0.008

0.015

0.006 0.01 0.004 0.005 0 0

0.002 0.0002

0.0006

0.001

0.0014

0.0018

0 0

0.0002

0.0004

0.0006

Figure 3.1: Error for Tikhonov’s regularization method.

0.0008

0.001

106

Regularization by Discretization 1

0.5

0.9

0.45

0.8

0.4

0.7

0.35

0.6

0.3

0.5

0.25

0.4

0.2

0.3

0.15

0.2

0.1

0.1

0.05

0 0

10

20

30

40

50

60

70

80

90

100

0.04

0 0

50

100

150

200

250

300

0.02 0.018

0.035

0.016 0.03 0.014 0.025

0.012

0.02

0.01 0.008

0.015

0.006 0.01 0.004 0.005 0 0

0.002 50

100

150

200

250

300

350

400

450

500

0

50

100

150

200

250

300

350

400

450

500

Figure 3.2: Error for Landweber’s method (a = 0.5).

1

0.5

0.9

0.45

0.8

0.4

0.7

0.35

0.6

0.3

0.5

0.25

0.4

0.2

0.3

0.15

0.2

0.1

0.1

0.05

0 2

3

4

5

6

7

8

9

10

0 0

2

4

6

8

10

12

14

16

18

20

0.02

0.04

0.018

0.035

0.016 0.03 0.014 0.025

0.012 0.01

0.02

0.008

0.015

0.006 0.01 0.004 0.005 0 0

0.002 5

10

15

20

25

30

0 5

6

7

8

Figure 3.3: Error for the conjugate gradient method.

9

10

3.5

Numerical Experiments for Symm’s Equation

107

In Figure 3.3, we show the results for the conjugate gradient method for the same example where again δ = 0.1, δ = 0.01, δ = 0.001, and δ = 0. The errors in the solution and the defects are again plotted versus the iteration number m. Here, we observe the same behavior as for Tikhonov’s method. We note the difference in the results for the Landweber method and the conjugate gradient method. The latter decreases the errors very quickly but is very sensitive to the exact stopping rule, while the Landweber iteration is slow but very stable with respect to the stopping parameter τ . The minimal values are errδ=0.1 ≈ 0.177, errδ=0.01 ≈ 0.0352, and errδ=0.001 ≈ 0.0054 for the Landweber iteration and errδ=0.1 ≈ 0.172, errδ=0.01 ≈ 0.0266, and errδ=0.001 ≈ 0.0038 for the conjugate gradient method. The corresponding factors are considerably larger than 102/3 ≈ 4.64 indicating the optimality of these methods also for smooth solutions (see the remarks following Theorem 2.15). Table 3.1: Least squares method n

δ = 0.1

δ = 0.01

1 2 3 4 5 6 10 12 15 20

38.190 15.772 5.2791 1.6209 1.0365 1.1954 2.7944 3.7602 4.9815 7.4111

38.190 15.769 5.2514 1.4562 0.3551 0.1571 0.2358 0.3561 0.4871 0.7270

δ = 0.001

δ=0

38.190 15.768 5.2511 1.4541 3.433 ∗ 10−1 7.190 ∗ 10−2 2.742 ∗ 10−2 3.187 ∗ 10−2 4.977 ∗ 10−2 7.300 ∗ 10−2

38.190 15.768 5.2511 1.4541 3.432 ∗ 10−1 7.045 ∗ 10−2 4.075 ∗ 10−5 5.713 ∗ 10−7 5.570 ∗ 10−10 3.530 ∗ 10−12

Table 3.2: Bubnov–Galerkin method n

δ = 0.1

δ = 0.01

1 2 3 4 5 6 10 12 15 20

38.190 15.771 5.2752 1.6868 1.1467 1.2516 2.6849 3.3431 4.9549 7.8845

38.190 15.769 5.2514 1.4565 0.3580 0.1493 0.2481 0.3642 0.4333 0.7512

δ = 0.001

δ=0

38.190 15.768 5.2511 1.4541 3.434 ∗ 10−1 7.168 ∗ 10−2 2.881 ∗ 10−2 3.652 ∗ 10−2 5.719 ∗ 10−2 7.452 ∗ 10−2

38.190 15.768 5.2511 1.4541 3.432 ∗ 10−1 7.045 ∗ 10−2 4.075 ∗ 10−5 5.713 ∗ 10−7 5.570 ∗ 10−10 3.519 ∗ 10−12

Next, we compute the same example using some projection methods. First, we list the results for the least squares method and the Bubnov–Galerkin method of Subsections 3.2.1 and 3.2.3 in Tables 3.1 and 3.2. We observe that both methods produce almost the same results, which reflect the estimates of Theorem 3.18. Note that for δ = 0 the error decreases exponentially with m. This

108

Regularization by Discretization

  reflects the fact that the best approximation min ψ − φn L2 : φn ∈ Xn converges to zero exponentially due to the analyticity of the solution ψ(s) = exp(3 sin s) (see [168], Theorem 11.7).

Now we turn to the collocation methods of Section 3.4. To implement the collocation method (3.83) for Symm’s integral equation and the basis functions (3.89), (3.93a), and (3.93b), we have to compute the integrals

1 − π

2π

eijs ln |γ(tk ) − γ(s)| ds ,

(3.101a)

0

j = −m, . . . , m − 1, k = 0, . . . , 2m − 1, and 1 − π

2π x ˆj (s) ln |γ(tk ) − γ(s)| ds ,

j, k = 0, . . . , 2m − 1 ,

(3.101b)

0

respectively. For the first integral (3.101a), we write using (3.64),

−

1 π

2π

eijs ln |γ(tk ) − γ(s)| ds

0

1 = − 2π

2π

ijs

e

% ln 4 sin

2 tk

0

= εj eijtk −

1 2π

2π

eijs ln

0

1 − s& ds − 2 2π

2π

eijs ln

0

|γ(tk ) − γ(s)|2 ds 4 sin2 (tk − s)/2

|γ(tk ) − γ(s)|2 ds , 4 sin2 (tk − s)/2

where εj = 0 for j = 0 and εj = 1/|j| otherwise. The remaining integral is computed by the trapezoidal rule. The computation of (3.101b) is more complicated. By Definition (3.93a), (3.93b) of x ˆj , we have to calculate tj +π/(2m) 

2

π/(2m) 

ln |γ(tk ) − γ(s)| ds = tj −π/(2m)

−π/(2m)

ln |γ(tk ) − γ(s + tj )|2 ds .

3.5

Numerical Experiments for Symm’s Equation

109

For j = k, the integrand is analytic, and we use Simpson’s rule π/(2m) 

g(s) ds ≈

n 

where π π − , s =

m n 2m

w g(s ) ,

=0

−π/(2m)

π w = · 3m n

'

1, 4, 2,

= 0 or n,

= 1, 3, . . . , n − 1,

= 2, 4, . . . , n − 2,

= 0, . . . , n. For j = k, the integral has a weak singularity at s = 0. We split the integrand into π/(2m) 

−π/(2m)

s& ds + ln 4 sin 2

π = −2 π/(2m)

%

π/(2m) 

2

ln −π/(2m) π/(2m) 

% s& ds + ln 4 sin2 2

|γ(tk ) − γ(s + tk )|2 ds 4 sin2 (s/2)

ln −π/(2m)

|γ(tk ) − γ(s + tk )|2 ds 4 sin2 (s/2)

  $π  because ln 4 sin2 (s/2) is even and 0 ln 4 sin2 (s/2) ds = 0 by (3.64). Both integrals are approximated by Simpson’s rule. For the same example as earlier, with 100 integration points for Simpson’s rule we obtain the following results for basis functions (3.89) (in Table 3.3) and basis functions (3.93a), (3.93b) (in Table 3.4). 

Table 3.3: Collocation method for basis functions (3.89) m

δ = 0.1

δ = 0.01

δ = 0.001

δ=0

1 2 3 4 5 6 10 12 15 20

6.7451 1.4133 0.3556 0.2525 0.3096 0.3404 0.5600 0.6974 0.8017 1.1539

6.7590 1.3877 2.791 ∗ 10−1 5.979 ∗ 10−2 3.103 ∗ 10−2 3.486 ∗ 10−2 5.782 ∗ 10−2 6.766 ∗ 10−2 8.371 ∗ 10−2 1.163 ∗ 10−1

6.7573 1.3880 2.770 ∗ 10−1 5.752 ∗ 10−2 1.110 ∗ 10−2 3.753 ∗ 10−3 5.783 ∗ 10−3 6.752 ∗ 10−3 8.586 ∗ 10−3 1.182 ∗ 10−2

6.7578 1.3879 2.769 ∗ 10−1 5.758 ∗ 10−2 1.099 ∗ 10−2 1.905 ∗ 10−3 6.885 ∗ 10−7 8.135 ∗ 10−9 6.436 ∗ 10−12 1.806 ∗ 10−13

110

Regularization by Discretization

Table 3.4: Collocation method for basis functions (3.93a) and (3.93b) m

δ = 0.1

δ = 0.01

δ = 0.001

δ=0

1 2 3 4 5 6 10 12 15 20 30

6.7461 1.3829 0.4944 0.3225 0.3373 0.3516 0.5558 0.6216 0.8664 1.0959 1.7121

6.7679 1.3562 4.874 ∗ 10−1 1.971 ∗ 10−1 1.649 ∗ 10−1 1.341 ∗ 10−1 8.386 ∗ 10−2 7.716 ∗ 10−2 9.091 ∗ 10−2 1.168 ∗ 10−1 1.688 ∗ 10−1

6.7626 1.3599 4.909 ∗ 10−1 2.000 ∗ 10−1 1.615 ∗ 10−1 1.291 ∗ 10−1 6.140 ∗ 10−2 4.516 ∗ 10−2 3.137 ∗ 10−2 2.121 ∗ 10−2 1.862 ∗ 10−2

6.7625 1.3600 4.906 ∗ 10−1 2.004 ∗ 10−1 1.617 ∗ 10−1 1.291 ∗ 10−1 6.107 ∗ 10−2 4.498 ∗ 10−2 3.044 ∗ 10−2 1.809 ∗ 10−2 8.669 ∗ 10−3

The difference for δ = 0 reflects the fact that the best approximation   min ψ − φn L2 : φn ∈ span {ˆ xj : j ∈ J} converges to zero exponentially for x ˆj defined by (3.89), while it converges to zero only of order 1/n for x ˆj defined by (3.93a) and (3.93b) (see Theorem 3.28). We have seen in this section that the theoretical investigations of the regularization strategies are confirmed by the numerical results for Symm’s integral equation.

3.6

The Backus–Gilbert Method

In this section, we study a different numerical method for “solving” finite moment problems of the following type: b kj (s) x(s) ds = yj ,

j = 1, . . . , n .

(3.102)

a

Here, yj ∈ R are any given numbers and kj ∈ L2 (a, b) arbitrary given functions. Certainly, we have in mind that yj = y(tj ) and kj = k(tj , ·). In Section 3.4, we studied the moment solution of such problems; see [184, 229]. We saw that the moment solution xn is a finite linear combination of the functions {k1 , . . . , kn }. Therefore, the moment solution xn is as smooth as the functions kj even if the true solution is smoother. The concept originally proposed by Backus and Gilbert ([13, 14]) does not primarily wish to solve the moment problem but rather wants to determine how well all possible models x can be recovered pointwise.

3.6

The Backus–Gilbert Method

111

Define the finite-dimensional operator K : L2 (a, b) → Rn by b (Kx)j =

kj (s) x(s) ds ,

j = 1, . . . , n,

x ∈ L2 (a, b) .

(3.103)

a

We try to find a left inverse S, that is, a linear operator S : Rn → L2 (a, b) such that (3.104) SKx ≈ x for all x ∈ L2 (a, b) . Therefore, SKx should be a simultaneous approximation to all possible x ∈ L2 (a, b). Of course, we have to make clear the meaning of the approximation. The general form of a linear operator S : Rn → L2 (a, b) has to be (Sy)(t) =

n 

t ∈ (a, b) ,

yj ϕj (t) ,

y = (yj ) ∈ Rn ,

(3.105)

j=1

for some ϕj ∈ L2 (a, b) that are to be determined from the requirement (3.104): (SKx)(t)

n 

=

b ϕj (t)

j=1

kj (s) x(s) ds a

⎡ ⎤ b  n ⎣ = kj (s) ϕj (t)⎦ x(s) ds . a

j=1

The requirement SKx ≈ x leads to the problem of approximating Dirac’s delta

n distribution δ(s − t) by linear combinations of the form j=1 kj (s) ϕj (t). For example, one can show that the minimum of 2 b b   n    ds dt k (s) ϕ (t) − δ(s − t) j j   a

a

j=1

(in the sense of distributions) is attained at ϕ(s) = A−1 k(s), where k(s) =   $b k1 (s), . . . , kn (s) and A ij = a ki (s) kj (s) ds, i, j = 1, . . . , n. For this minn imization criterion, x = j=1 yj ϕj is again the moment solution of Subsec−s -norm tion 3.4.1. In [184], it is shown that minimizing with respect to an Hper s -spaces. We refer also to [272] for s > 1/2 leads to projection methods in Hper for a comparison of several minimization criteria. The Backus–Gilbert method is based on a pointwise minimization criterion: Treat t ∈ [a, b] as a fixed parameter and determine the numbers ϕj = ϕj (t) for j = 1, . . . , n, as the solution of the following minimization problem: b minimize a

 2  n   |s − t|  kj (s) ϕj  ds 2

j=1

(3.106a)

112

Regularization by Discretization

subject to ϕ ∈ Rn and

b  n a

kj (s) ϕj ds = 1 .

(3.106b)

j=1

Using the matrix-vector notation, we rewrite this problem in short form: minimize

ϕ Q(t) ϕ

r · ϕ = 1,

subject to

where b Q(t)ij

=

|s − t|2 ki (s) kj (s) ds,

i, j = 1, . . . , n,

a

b rj

=

kj (s) ds ,

j = 1, . . . , n .

a

This is a quadratic minimization problem with one linear equality constraint. We assume that r = 0 because otherwise the constraint (3.106b) cannot be satisfied. Uniqueness and existence are assured by the following theorem, which also gives a characterization by the Lagrange multiplier rule. Theorem 3.29 Assume that {k1 , . . . , kn } are linearly independent. Then the symmetric matrix Q(t) ∈ Rn×n is positive definite for every t ∈ [a, b]. The minimization problem (3.106a), (3.106b) is uniquely solvable. ϕ ∈ Rn is a solution of (3.106a) and (3.106b) if and only if there exists a number λ ∈ R (the Lagrange multiplier) such that (ϕ, λ) ∈ Rn × R solves the linear system Q(t)ϕ − λ r = 0

and

r · ϕ = 1.

(3.107)

λ = ϕ Q(t) ϕ is the minimal value of this problem. Proof:

From 

b

ϕ Q(t) ϕ =

 n 2    |s − t|  kj (s) ϕj  ds, 2

a

j=1

we conclude first that ϕ Q(t) ϕ ≥ 0 and second that ϕ Q(t) ϕ = 0 implies

n k that j=1 j (s) ϕj = 0 for almost all s ∈ (a, b). Because {kj } are linearly independent, ϕj = 0 for all j follows. Therefore, Q(t) is positive definite. Existence, uniqueness, and equivalence to (3.107) are elementary results from optimization theory; see [269].  n  Definition 3.30 We denote by ϕj (t) j=1 ∈ Rn the unique solution ϕ ∈ Rn of (3.106a) and (3.106b). The Backus–Gilbert solution xn of b kj (s) xn (s) ds = yj , a

j = 1, . . . , n ,

3.6

The Backus–Gilbert Method

is defined as xn (t) =

n 

113

yj ϕj (t) ,

t ∈ [a, b] .

(3.108)

j=1

The minimal value λ = λ(t) = ϕ(t) Q(t)ϕ(t) is called the spread.

n We remark that, in general, the Backus–Gilbert solution xn = j=1 yj ϕj $b is not a solution of the moment problem, that is, a kj (s) xn (s) ds = yj ! This is certainly a disadvantage. On the other hand, the solution x is analytic in [a, b]—even for nonsmooth data kj . We can prove the following lemma. Lemma 3.31 ϕj and λ are rational functions. More precisely, there exist polynomials pj , q ∈ P2(n−1) and ρ ∈ P2n such that ϕj = pj /q, j = 1, . . . , n, and λ = ρ/q. The polynomial q has no zeros in [a, b]. Proof: Obviously, Q(t) = Q0 − 2t Q1 + t2 Q2 with symmetric matrices Q0 , n Q1 , Q2 . We search for a polynomial solution p ∈ Pm and ρ ∈ Pm+2 of Q(t)p(t) − ρ(t) r = 0 with m = 2(n − 1). Because the number of equations is n(m + 3) = 2n2 + n and the number of unknowns is n(m n + 1) + (m + 3) = 2n2 + n + 1, there exists a nontrivial solution p ∈ Pm and ρ ∈ Pm+2 . If p(tˆ) = 0 for some tˆ ∈ [a, b], then ρ(tˆ) = 0 because r = 0. In this case, we divide the equation by (t − tˆ). Therefore, we can assume that p has no zero in [a, b]. Now we define q(t) := r · p(t) for t ∈ [a, b]. Then q ∈ Pm has no zero in [a, b] because otherwise we would have 0 = ρ(tˆ) r · p(tˆ) = p(tˆ) Q(tˆ) p(tˆ) ; thus p(tˆ) = 0, a contradiction. Therefore, ϕ := p/q and λ := ρ/q solves (3.107). By the uniqueness result, this is the only solution.  For the following error estimates, we assume two kinds of a priori information on x depending on the norm of the desired error estimate. Let Xn = span{kj : j = 1, . . . , n} . 2 Theorem 3.32 Let

nx ∈ L (a, b) be any solution of the finite moment problem (3.102) and xn = j=1 yj ϕj be the Backus–Gilbert solution. Then the following error estimates hold:

(a) Assume that x is Lipschitz continuous with constant > 0, that is, |x(t) − x(s)| ≤ |s − t| Then

for all s, t ∈ [a, b] .

√ |xn (t) − x(t)| ≤ b − a n (t)

(3.109)

for all n ∈ N, t ∈ [a, b], where n (t) is defined by b  b |s−t|2 |zn (s)|2 ds : zn ∈ Xn , zn (s) ds = 1 . (3.110) 2n (t) := min a

a

114

Regularization by Discretization

(b) Let x ∈ H 1 (a, b). Then there exists c > 0, independent of x, such that xn − xL2 ≤ c x L2 n ∞

for all n ∈ N .

(3.111)

Proof: By the definition of the Backus–Gilbert solution and the constraint on ϕ, we have n 

xn (t) − x(t) =

yj ϕj (t) − x(t)

j=1

a

n b 

=

b  n

kj (s) ϕj (t) ds

j=1

  kj (s) x(s) − x(t) ϕj (t) ds .

j=1 a

Thus

 b   n    |x(s) − x(t)| ds . k (s) ϕ (t) j j  

|xn (t) − x(t)| ≤

a

j=1

Now we distinguish between parts (a) and (b): (a) Let |x(t) − x(s)| ≤ |t − s|. Then, by the Cauchy–Schwarz inequality and the definition of ϕj , b |xn (t) − x(t)|

≤

a

≤



√

   n  1 ·  kj (s) ϕj (t) |t − s| ds j=1

⎡ b ⎤1/2 2    n  b − a ⎣  kj (s) ϕj (t) |t − s|2 ds⎦ a

=

√

j=1

b − a n (t) .

(b) First, we define the cutoff function λδ on [a, b] × [a, b] by  1, |t − s| ≥ δ, λδ (t, s) = 0, |t − s| < δ. Then, by the Cauchy–Schwarz inequality again, ⎤2 ⎡ b     n  ⎣ λδ (t, s)  kj (s) ϕj (t) |x(s) − x(t)| ds⎦  j=1

a

⎤2 ⎡ b       n   x(s) − x(t)   ds⎦ = ⎣  kj (s) ϕj (t) (t − s) λδ (t, s)  t−s  a

≤

j=1 2

b

n (t)

a

   x(s) − x(t) 2  ds . λδ (t, s)  s−t 

(3.112)

3.6

The Backus–Gilbert Method

115

Integration with respect to t yields ⎤2 ⎡   b b  n  ⎣ λδ (t, s)  kj (s) ϕj (t) |x(s) − x(t)| ds⎦ dt  a

≤

j=1

a

n 2∞

  x(s) − x(t) 2    s − t  λδ (t, s) ds dt .

b

b 

a

a

The following technical lemma from the theory of Sobolev spaces yields the assertion.  Lemma 3.33 There exists c > 0 such that  b b   x(s) − x(t) 2  2    s − t  λδ (t, s) ds dt ≤ cx L2 a

a

for all δ > 0 and x ∈ H 1 (a, b). Here, the cutoff function λδ is defined by (3.112). Proof:

First, we estimate |x(s) − x(t)|

2

 t 2  t       =  1 · x (τ ) dτ  ≤ |t − s|  |x (τ )|2 dτ  s

s

and thus, for s = t,    t   x(s) − x(t) 2   1   ≤  |x (τ )|2 dτ  .  s−t    |s − t| s

Now we fix t ∈ (a, b) and write  b   x(s) − x(t) 2    s − t  λδ (t, s) ds a

t ≤ a

b =

λδ (t, s) t−s

t



2

b

|x (τ )| dτ ds + s

t

λδ (t, s) s−t

s

|x (τ )|2 dτ ds

t

|x (s)|2 Aδ (t, s) ds ,

a

where

Aδ (t, s) =

⎧ $s ⎪ ⎪ ⎨ ⎪ ⎪ ⎩

a $b s

λδ (t,τ ) |t−τ | λδ (t,τ ) |t−τ |

dτ,

'

a ≤ s < t, =

dτ,

t < s ≤ b,

t−a ln max(δ,t−s) , s ≤ t, b−t , s ≥ t. ln max(δ,s−t)

116

Regularization by Discretization

Finally, we estimate b b a

a

2

|x(s) − x(t)| λδ (t, s) ds dt ≤ |s − t|2

and

b

⎛ b ⎞  |x (s)|2 ⎝ Aδ (t, s) dt⎠ ds

a

a

b Aδ (t, s) dt ≤ c for all s ∈ (a, b) and δ > 0 a

which is seen by elementary integration.



From these error estimates, we observe that the rate of convergence depends on the magnitude of n , that is, how well the kernels approximate the delta distribution. Finally, we study the question of convergence for n → ∞. Theorem 3.34 Assume that {kj : j ∈ N} is linearly independent and dense in L2 (a, b). Then n ∞ −→ 0 for n → ∞ .   Proof: For fixed t ∈ [a, b] and arbitrary δ ∈ 0, (b − a)/2 , we define ⎤−1 ⎡ b  1  , |s − t| ≥ δ, |s−t| v˜(s) := and v(s) := ⎣ v˜(τ ) dτ ⎦ v˜(s) . 0, |s − t| < δ, $b

a



Then v ∈ L2 (a, b) and a v(s) ds = 1. Because Xn is dense in L2 (a, b), there exists a sequence v˜n ∈ Xn with v˜n → v in L2 (a, b). This implies also that $b $b v˜ (s) ds → a v(s) ds = 1. Therefore, the functions a n ( b )−1 vn := v˜n (s) ds v˜n ∈ Xn a

converge to v in L2 (a, b) and are normalized by admissible, and we conclude that 2

n (t)

b ≤

$b a

vn (s) ds = 1. Thus vn is

|s − t|2 vn (s)2 ds

a

b

2

2

b

|s − t| v(s) ds + 2

= a

  |s − t|2 v(s) vn (s) − v(s) ds

a

b +

 2 |s − t|2 vn (s) − v(s) ds

a

≤

"b #−2   (b − a) v˜(s) ds + (b − a)2 2 vL2 vn − vL2 + vn − v2L2 . a

3.7

Problems

117

This shows that lim sup n (t) ≤

√

n→∞

"b b−a

#−1 v˜(s) ds

for all t ∈ [a, b] .

a

Direct computation yields b v˜(s) ds ≥ c + | ln δ| a

for some c independent of δ; thus √ b−a lim sup n (t) ≤ c + | ln δ| n→∞

  for all δ ∈ 0, (b − a)/2 .

This yields pointwise convergence, that is, n (t) → 0 (n → ∞) for every t ∈ [a, b]. Because n (t) is monotonic with respect to n, Dini’s well-known theorem from classical analysis (see, for example, [231]) yields uniform convergence.  For further aspects of the Backus–Gilbert method, we refer to [32, 114, 144, 162, 163, 244, 272, 273].

3.7

Problems

3.1 Let Qn : C[a, b] → S1 (t1 , . . . , tn ) be the interpolation operator from Example 3.3. Prove that Qn L(C[a,b]) = 1 and derive an estimate of the form Qn x − x∞ ≤ c h x ∞ for x ∈ C 1 [a, b], where h = max{tj − tj−1 : t = 2, . . . , n}. 3.2 Let K : X → X be self-adjoint and positive definite and let y ∈ X. Define ψ(x) = (Kx, x)X − 2 Re(y, x)X for x ∈ X. Prove that x∗ ∈ X is a minimum of ψ if and only if x∗ solves Kx∗ = y. 3.3 Define the space Xn by Xn =



 aj eijt : aj ∈ C

|j|≤n

and let Pn : L2 (0, 2π) → Xn be the orthogonal projection operator. Prove that for r ≥ s there exists c > 0 such that r ψn Hper

≤

s Pn ψ − ψHper

≤

s c nr−s ψn Hper for all ψn ∈ Xn , 1 r r c r−s ψHper for all ψ ∈ Hper (0, 2π) . n

118

Regularization by Discretization

3.4 Show that the worst-case error of Symm’s equation under the information s ψHper ≤ E for some s > 0 is given by   s ≤ c δ s/(s+1) . F δ, E,  · Hper 3.5 Let Ω ⊂ R2 be the disk of radius a = exp(−1/2). Then ψ = 1 is the unique solution of Symm’s integral equation (3.57) for f = 1. Compute explicitly the errors of the least squares solution, the dual least squares solution, and the Bubnov–Galerkin solution as in Section 3.3, and verify that the error estimates of Theorem 3.18 are asymptotically sharp. 3.6 Let tk = k/n, k = 1, . . . , n, be equidistant collocation points. Let Xn be the space of piecewise constant functions as in (3.81) andPn : L2 (0, 1) → Xn be the orthogonal projection operator. Prove that Xn is dense in L2 (0, 1) and 1  x − Pn xL2 ≤ x L2 n for all x ∈ H 1 (0, 1) (see Problem 3.1). 3.7 Show that the moment solution can also be interpreted as the solution of a dual least squares method. 3.8 Consider moment collocation of the equation t x(s) ds = y(t),

t ∈ [0, 1] ,

0

in the space Xn = S1 (t1 , . . . , tn ) of linear splines. Show that the moment solution xn coincides with the two-sided difference quotient, that is, xn (tj ) =

   1   y tj+1 + h − y tj−1 − h , 2h

where h = 1/n. Derive an error estimate for xδn −xL2 as in Example 3.23.

Chapter 4

Nonlinear Inverse Problems In the previous chapters, we considered linear problems which we wrote as Kx = y, where K was a linear and (often) compact operator between Hilbert spaces. Needless to say that most problems in applications are nonlinear. For example, even in the case of a linear differential equation of the form −u + cu = f for the function u the dependence of u on the parameter function c is nonlinear; that is, the mapping c → u is nonlinear.1 In Chapters 5, 6, and 7 we will study particular nonlinear problems to determine parameters of an ordinary or partial differential equation from the knowledge of the solution. Although we believe that the best strategies for solving nonlinear problems are intrinsically linked to the particular nature of the underlying problem, there are general methods for solving these nonlinear problems if they can be written in the form K(x) = y, where K is now a nonlinear operator between Hilbert spaces or Banach spaces. Guided by the structure of Chapter 2, we will study the nonlinear form of the Tikhonov regularization in Section 4.2 and the extension of the Landweber method in Section 4.3. Since the investigation of the latter one is already rather complicated, we do not present the extension of the conjugate gradient method or methods of Newton type, but refer the interested reader to the monograph [149] of Kaltenbacher, Neubauer, and Scherzer. We start this chapter with a clarification of the notion of ill-posedness in Section 4.1 and its relation to the ill-posedness of the linearized problem. In Section 4.2, we study Tikhonov’s regularization method for nonlinear problems. In contrast to the linear case, the question of existence of a global minimum of the Tikhonov functional is not obvious and requires more advanced tools from functional analysis, in particular, on weak topologies. We include the arguments in Subsection 4.2.1 but note already here that this Subsection is not central for the understanding of the further theory. It can be omitted because we formulate the existence of a minimum also as Assumption 4.9 at the beginning of Subsection 4.2.2. After the application of the general theory to the abovementioned parameter identification problem for a boundary value

1 The differential equation has to be complemented by initial or boundary conditions, of course.

© Springer Nature Switzerland AG 2021 A. Kirsch, An Introduction to the Mathematical Theory of Inverse Problems, Applied Mathematical Sciences 120, https://doi.org/10.1007/978-3-030-63343-1 4

119

120

Nonlinear Inverse Problems

problem for an ordinary differential equation, in Subsection 4.2.4, we present some of the basic ideas for Tikhonov’s method in Banach spaces and more general metrics to penalize the discrepancy and to measure the error. As a particular example, we consider the determination of a sparse solution of a linear equation. Further, tools from convex analysis are needed, such as the subdifferential and the Bregman distance. Finally, in Section 4.3, we return to the Hilbert space setting and extend the Landweber method from Sections 2.3 and 2.6 to the nonlinear case.

4.1

Local Illposedness

In this chapter, we assume that X and Y are normed spaces (in most cases Hilbert spaces), K : X ⊃ D(K) → Y a nonlinear mapping with domain of definition D(K) ⊂ X. Let x∗ ∈ D(K) and y ∗ ∈ Y such that K(x∗ ) = y ∗ . It is the aim—as in the linear case—to determine an approximate solution to x∗ when the right-hand side y ∗ is perturbed; that is, replaced by y δ ∈ Y with y δ − y ∗ Y ≤ δ. In the following, let B(x, r) = {z : x − zX < r} denote the open ball centered at x with radius r. The following notion of local ill-posedness goes back to Hofmann and Scherzer (see, e.g., [137]). Definition 4.1 Let x∗ ∈ D(K) and y ∗ = K(x∗ ). The equation K(x) = y is called locally improperly-posed or locally ill-posed at x∗ if for any sufficiently small ρ > 0 there exists a sequence xn ∈ D(K) ∩ B(x∗ , ρ) such that K(xn ) → K(x∗ ) but (xn ) does not converge to x∗ as n tends to infinity. Example 4.2   1  1] × [0, 1] × R , k = k(t, s, r), and let there exist c1 > 0 with Let k ∈ C [0, ∂k(t, s, r)/∂r ≤ c1 for all (t, s, r) ∈ [0, 1] × [0, 1] × R. Define 1 K(x)(t) =

  k t, s, x(s) ds ,

t ∈ [0, 1] ,

for x ∈ L2 (0, 1) .

0

Then K is well-defined from L2 (0, 1) into itself, and the equation K(x) = y is locally ill-posed at x = 0. Proof: Let x ∈ L2 (0, 1). The application of the fundamental theorem of calculus in the form r ∂k (t, s, r ) dr , r ∈ R , k(t, s, r) = k(t, s, 0) + (4.1) ∂r 0

 2 2      implies k(t, s, r) ≤ k(t, s, 0) + c1 |r|, thus k(t, s, r) ≤ 2k(t, s, 0) + 2c21 r2 , thus, using the Cauchy-Schwarz inequality,   K(x)(t)2 ≤

1 0

   k t, s, x(s) 2 ds ≤ 2

 0

  k(t, s, 0)2 + c21 x(s)2 ds

1 

4.1

Local Illposedness

121

2  for all t ∈ [0, 1]. Therefore, K(x) is measurable and K(x)(t) is bounded which implies K(x) ∈ L2 (0, 1). √ Let now ρ > 0 be arbitrary and xn (t) = ρ 2n + 1 tn , t ∈ [0, 1]. Then xn 2L2 (0,1) = 1 ρ2 (2n + 1) t2n dt = ρ2 and, with (4.1), 0

  K(xn )(t) − K(0)(t)

1 ≤

c1

  √ xn (s) ds = c1 ρ 2n + 1

0

=

1

sn ds

0

√ c1 ρ 2n + 1 → 0 as n → ∞ . n+1

Therefore, the equation K(x) = y is locally improperly-posed at x = 0. A second, and more concrete, example is formulated as Problem 4.2. If K is continuously Fr´echet-differentiable at x∗ with Lipschitz continuous derivative then local illposedness of the nonlinear problem implies the illposedness of the linearization. Theorem 4.3 Let K be Fr´echet-differentiable in the ball B(x∗ , ρˆ) and let there exists γ > 0 with K  (x)−K  (x∗ )L(X,Y ) ≤ γx−x∗ X for all x ∈ B(x∗ , ρˆ). Let the equation K(x) = y be locally ill-posed at x∗ . Then K  (x∗ ) is not boundedly invertible; that is, the linear equation K  (x∗ )h = z is also ill-posed. Proof: We assume on the contrary that K  (x∗ ) is boundedly invertible and  ∗ −1 L(Y,X) =: q < 1. For this ρ we choose choose ρ ∈ (0, ρˆ) such that ργ 2 K (x ) a sequence xn ∈ B(x∗ , ρ) by Definition 4.1. Lemma A.63 of the Appendix A.7 implies the representation K(xn ) − K(x∗ ) = K  (x∗ )(xn − x∗ ) + rn

with

rn Y ≤

for all n ∈ N; that is,   K  (x∗ )−1 K(xn ) − K(x∗ ) = xn − x∗ + K  (x∗ )−1 rn

γ xn − x∗ 2X 2

for all n ∈ N ,

thus xn − x∗ X

≤ ≤

≤

 ∗ −1   K (x ) K(xn ) − K(x∗ ) + K  (x∗ )−1 rn X X

K  (x∗ )−1 L(Y,X) K(xn ) − K(x∗ )Y γ + K  (x∗ )−1 L(Y,X) xn − x∗ 2X 2 K  (x∗ )−1 L(Y,X) K(xn ) − K(x∗ )Y + q xn − x∗ X ,

thus (1 − q) xn − x∗ X ≤ K  (x∗ )−1 L(Y,X) K(xn ) − K(x∗ )Y ,

122

Nonlinear Inverse Problems

and this expression converges to zero because of K(xn ) → K(x∗ ). This contra dicts (xn ) does not converge to x∗ . We will see in Section 4.3 that also the reverse assertion is true provided an addition condition (“tangential cone condition”) is satisfied, see Lemma 4.35.

4.2

The Nonlinear Tikhonov Regularization

In this section (except of Subsection 4.2.4), we assume that X and Y are Hilbert spaces with inner products (·, ·)X and (·, ·)Y , respectively, and corresponding norms. Let again x∗ ∈ D(K) and y ∗ ∈ Y with K(x∗ ) = y ∗ and y δ ∈ Y with y δ − y ∗ Y ≤ δ. Let, in addition, x ˆ ∈ X be given which is thought of being an approximation of the true solution x∗ . We define the Tikhonov functional by ˆ2X , Jα,δ (x) = K(x) − y δ 2Y + αx − x

x ∈ D(K) .

(4.2)

In the first subsection, we will discuss briefly the question of existence of minima of the Tikhonov functional and stability with respect to perturbation of y δ . This part needs some knowledge of the weak topology in Hilbert spaces. We have collected the needed results in Section A.9 of the Appendix for the convenience of the reader. One can easily drop this subsection if one is not familiar (or interested) with this part of functional analysis. The further analysis is quite independent of this subsection.

4.2.1

Existence of Solutions and Stability

We recall the following definition (see remark following Definition A.75 of the Appendix A.9). Definition 4.4 Let X be a Hilbert space. A sequence (xn ) in X is said to converge weakly to x ∈ X if limn→∞ (xn , z)X = (x, z)X for all z ∈ X. If Y is a second Hilbert space and K : X ⊃ D(K) → Y is a (nonlinear) mapping then K is called weak-to-weak continuous, if K maps weakly convergent sequences in D(K) into weakly convergent sequences in Y . We will need the following two results (see Corollary A.78 and part (d) of Theorem A.76). Theorem 4.5 Let X be a Hilbert space. (a) Every bounded sequence (xn ) in X contains a weak accumulation point; that is, a weakly convergent subsequence. (b) Every convex and closed set U ⊂ X is also weakly closed; that is, if the sequence (xn ) in U converges weakly to some x ∈ X then necessarily x ∈ U.

4.2

The Nonlinear Tikhonov Regularization

123

The norm function  · X fails to be weakly continuous but has the following property which is sometimes called “weak lower semi-continuity”. Lemma 4.6 Let X be a Hilbert space and let the sequence (xn ) converge weakly to x. Then lim inf xn − zX ≥ x − zX n→∞

Proof:

for all z ∈ X .

This follows from the formula xn − z2X − x − z2X

=

2 Re(xn − x, x − z)X + xn − x2X

≥

2 Re(xn − x, x − z)X

because the right hand side of the inequality tends to zero as n tends to infinity.  Now we are able to prove the existence of minima of the Tikhonov functional under appropriate assumptions. Theorem 4.7 Let X and Y be Hilbert spaces, y δ ∈ Y , and K : X ⊃ D(K) → Y weak-to-weak continuous with convex and closed domain of definition D(K). Then there exists a global minimum of Jα,δ , defined in (4.2), on D(K) for all α > 0. Proof: Let xn ∈ D(K) be a minimizing sequence; that is, Jα,δ (xn ) → J ∗ := inf{Jα,δ (x) : x ∈ D(K)} as n → ∞. From αxn − x ˆ2X ≤ Jα,δ (xn ) ≤ J ∗ + 1 for sufficiently large n, we conclude that the sequence (xn ) is bounded. By part (a) of Theorem 4.5, there exists a subsequence—which we also denote by (xn )— which converges weakly to some x. Also, x ∈ D(K) by part (b) of this theorem. Furthermore, K(xn ) converges weakly to K(x) by the assumption on ˆX ≥ x − x ˆX and lim inf K(xn ) − K. Lemma 4.6 implies that lim inf xn − x n→∞

n→∞

y δ Y ≥ K(x) − y δ Y . Therefore, for any ε > 0 there exists N ∈ N such that for all n ≥ N Jα,δ (x)

= αx − x ˆ2X + K(x) − y δ 2Y ≤ ≤

αxn − x ˆ2X + K(xn ) − y δ 2Y + ε = Jα,δ (xn ) + ε J ∗ + 2ε .

This holds for all ε > 0, thus J ∗ ≤ Jα,δ (x) ≤ J ∗ which proves that x is a  minimum of Jα,δ . In the same way one proves stability.

124

Nonlinear Inverse Problems

Theorem 4.8 Let the assumptions of the previous theorem hold, and in addition, (yn ) be a sequence with yn → y δ as n → ∞. Let xn ∈ D(K) be a minimum ˆ2X + K(x) − yn 2Y on D(K). Then there exist weak accuof Jn (x) := αx − x mulation points of the sequence (xn ), and every weak accumulation point x is a minimum of Jα,δ . If in addition, K is weak-to-norm continuous, then every weak accumulation point x of (xn ) is also an accumulation point with respect to the norm. Proof:

The estimate

αxn − x ˆ2X + K(xn ) − yn 2Y

= Jn (xn ) ≤ Jn (ˆ x) = K(ˆ x) − yn 2Y 2  ≤ K(ˆ x) − y δ Y + y δ − yn Y  2 ≤ K(ˆ x) − y δ Y + 1

for sufficiently large n implies again that the sequence (xn ) is bounded. Thus, it contains weak accumulation points by Theorem 4.5. Let x ∈ D(K) be a weak accumulation point, without loss of generality let xn itself converge weakly to x. Then K(xn ) − yn converges weakly to K(x) − y δ and thus as before ˆX ≥ x − x ˆX and lim inf K(xn ) − yn Y ≥ K(x) − y δ Y . Set lim inf xn − x n→∞

n→∞

Jn∗ = Jn (xn ). Then, for any ε > 0, there exists N ∈ N such that for all n ≥ N and all x ∈ D(K) Jα,δ (x)

≤ ≤

αxn − x ˆ2X + K(xn ) − yn 2Y + ε = Jn∗ + ε Jn (x) + ε = αx − x ˆ2X + K(x) − yn 2Y + ε .

(4.3)

Now we let n tend to infinity. This yields Jα,δ (x) ≤ lim inf Jn∗ + ε ≤ lim inf Jn∗ + ε ≤ Jα,δ (x) + ε . n→∞

n→∞

Letting ε tend to zero proves the optimality of x and also the convergence of Jn∗ to Jα,δ (x). Let now K be weak-to-norm continuous. From   ˆ2X − x − x ˆ2X = Jn∗ − Jα,δ (x) + K(x) − y δ 2Y − K(xn ) − yn 2Y α xn − x and the convergence of K(xn ) to K(x), we conclude that the right-hand side ˆX → x− x ˆX . Finally, the binomial formula converges to zero and thus xn − x yields xn − x2X which tends to zero.

4.2.2

= (xn − x ˆ) − (x − x ˆ)2X = xn − x ˆ2X + x − x ˆ2X − 2 Re(xn − x ˆ, x − x ˆ)X 

Source Conditions And Convergence Rates

We make the following general assumption.

4.2

The Nonlinear Tikhonov Regularization

125

Assumption 4.9 (a) D(K) is open, and x∗ ∈ D(K) ∩ B(ˆ x, ρ) is a solux, ρ), and K is continuously Fr´echettion of K(x) = y ∗ in some ball B(ˆ differentiable on D(K) ∩ B(ˆ x, ρ), (b) K  (x∗ ) is compact from X into Y . (c) The Tikhonov functional Jα,δ possesses global minima xα,δ ∈ D(K) on D(K) for all α, δ > 0; that is, ˆ2X + K(xα,δ ) − y δ 2Y ≤ αx − x ˆ2X + K(x) − y δ 2Y (4.4) αxα,δ − x for all x ∈ D(K). We refer to the previous subsection where we showed part (c) of this assumption under appropriate smoothness assumptions on K. Now we are able to formulate the condition which corresponds to the “source condition” in the linear case. We will include the more general form with index functions and introduce the following notion. Definition 4.10 Any monotonically increasing and continuous function ϕ : [0, δmax ] → R (for some δmax > 0) with ϕ(0) = 0 is called an index function. The most prominent examples for index functions are ϕ(t) = tσ for any σ > 0 but also ϕ(t) = −1/ ln t for 0 < t < 1 (and ϕ(0) = 0). By calculating the second derivative, one observes that the latter one is concave on [0, 1/e2 ] and the first class is concave whenever σ ≤ 1. The linear case ϕ(t) = βt for t ≥ 0 is particularly important. Assumption 4.11 (Source condition) Let D(K) be open and x∗ ∈ D(K) ∩ x, ρ) and let K be differenB(ˆ x, ρ) be a solution of K(x) = y ∗ in some ball B(ˆ tiable on D(K) ∩ B(ˆ x, ρ). Furthermore, let ϕ : [0, δmax ] → R be a concave index function with δmax ≥ K  (x∗ )L(X,Y ) . (i) Let K  be locally Lipschitz continuous; that is, there exists γ > 0 with K  (x) − K  (x∗ )L(X,Y ) ≤ γx − x∗ X

for all x ∈ B(x∗ , ρ) ∩ D(K) ,

(ii) and there exists w ∈ X with   ˆ = ϕ [(K  (x∗ )∗ (K  (x∗ )]1/2 w x∗ − x

and

γ wX < 1 .

We note that for a linear compact operator A : X → Y (in the present   case A := K  (x∗ )) the operator ϕ [A∗ A]1/2 from X into itself is defined as in (A.47) by a singular system {μj , xj , yj : j ∈ J} for A, see Appendix A.6, Theorem A.57, where J is finite or J = N, namely,

  ϕ [A∗ A]1/2 z = ϕ(μj ) (z, xj )X xj , z ∈ X . j∈J

126

Nonlinear Inverse Problems

 In the special case that ϕ(t) = tσ the condition reads as x∗ − x ˆ = (A∗ A)σ/2 w which is just the source condition of the linear case ˆ = 0). Again, for  (for x  σ = 1 the ranges R [(K  (x∗ )∗ (K  (x∗ )]1/2 and R K  (x∗ )∗ coincide which is seen from the singular system {μj , xj , yj : j ∈ J}. Therefore, in this linear case ˆ = K  (x∗ )∗ v for some v ∈ Y with γ vY < 1. part (ii) takes the form x∗ − x In the past decade, a different kind of source conditions has been developed which does not need the derivative of K. It can be generalized to a wider class of Tikhonov functionals with non-differentiable K acting between Banach spaces. We refer to Subsection 4.2.4 for a short glimpse on these extensions. Assumption 4.12 (Variational Source condition) Let D(K) be open and x∗ ∈ x, ρ) and ϕ : D(K) ∩ B(ˆ x, ρ) be a solution of K(x) = y ∗ in some ball  B(ˆ ∗ index function with δ ≥ sup K(x ) − K(x)Y : [0, δmax ] → R be a concave max x ∈ D(K) ∩ B(ˆ x, ρ) . Furthermore, there exists a constant β > 0 such that   ˆ2X − x∗ − x ˆ2X + ϕ K(x∗ ) − K(x)Y βx∗ − x2X ≤ x − x for all x ∈ B(ˆ x, ρ) ∩ D(K). This assumption is also known as a variational inequality (see, e.g., [245]). We prefer the notion of variational source condition as, e.g., in [96], because it takes the role of the source condition. We now show a relationship between Assumption 4.11 and Assumption 4.12. Theorem 4.13 Let D(K) be open and x∗ ∈ D(K) ∩ B(ˆ x, ρ) be a solution of K(x) = y ∗ with ρ ∈ (0, 1/2) and ϕ : [0, δmax ] → R be a concave index function  ∗ ∗ with δmax ≥ K (x )L(X,Y ) and δmax ≥ sup K(x ) − K(x)Y : x ∈ D(K) ∩ B(ˆ x, ρ) . (a) The variational source condition of Assumption 4.12 is equivalent to the following condition: There exists 0 ≤ σ < 1 such that   2 Re(x∗ − x ˆ, x∗ − x)X ≤ σ x∗ − x2X + ϕ K(x∗ ) − K(x)Y (4.5) for all x ∈ B(ˆ x, ρ) ∩ D(K). (b) Let Assumption 4.11 hold  and, in addition, ϕ(t) = t for all t or x, ρ) ∩ D(K) K  (x∗ )(x − x∗ )Y ≤ η K(x) − K(x∗ )Y for all x ∈ B(ˆ where η is another concave index function. Then also Assumption 4.12 holds with some index function ϕ˜ which is linear if ϕ is linear. (c) Let Assumption 4.12 hold for ϕ(t) = βt for all t ≥ 0. Then x∗ − x ˆ ∈   ∗ ∗   ∗ ∗  ∗ 1/2  R K (x ) = R K (x ) K (x ) . Proof:

(a) From the elementary equation x − x ˆ2X − x∗ − x ˆ2X = 2 Re(x∗ − x ˆ, x − x∗ )X + x − x∗ 2X

4.2

The Nonlinear Tikhonov Regularization

127

we observe that the variational source condition is equivalent to   βx − x∗ 2X ≤ 2 Re(x∗ − x ˆ, x − x∗ )X + x − x∗ 2X + ϕ K(x∗ ) − K(x)Y , that is,   2 Re(x∗ − x ˆ, x∗ − x)X ≤ (1 − β)x − x∗ 2X + ϕ K(x∗ ) − K(x)Y , which has the desired form with σ = 1 − β if β ≤ 1 and σ = 0 if β > 1. (b) Set A := K  (x∗ ) for abbreviation. Then, for x ∈ B(ˆ x, ρ) ∩ D(K),     Re(x∗ − x ˆ, x∗ − x)X = Re ϕ [A∗ A]1/2 w, x∗ − x X     = Re w, ϕ [A∗ A]1/2 (x∗ − x) X   ≤ wX ϕ [A∗ A]1/2 (x∗ − x) . X

Now we use the fact that for any concave index function  ∗ 1/2    ϕ [A A] z X ≤ ϕ AzY for all z ∈ X with zX ≤ 1 . For a proof, we refer to Lemma A.73 of the Appendix. Therefore,   Re(x∗ − x (4.6) ˆ, x∗ − x)X ≤ wX ϕ K  (x∗ )(x∗ − x)Y .   If K  (x∗ )(x − x∗ )Y ≤ η K(x) − K(x∗ )Y for all x ∈ B(ˆ x, ρ) ∩ D(K) then   Re(x∗ − x ˆ, x∗ − x)X ≤ wX ϕ η(K(x) − K(x∗ )Y ) because of the monotonicity of ϕ. This proves the estimate (4.5) with σ = 0 and ϕ˜ = 2wX ϕ◦η. Note that the composition of two concave index functions is again a concave index function. If ϕ(t) = t for all t we use the estimate K(x) − K(x∗ ) − K  (x∗ )(x − x∗ )Y ≤

γ ∗ x − x2X 2

for all x ∈ B(ˆ x, ρ) ∩ D(K) (see Lemma A.63 of the Appendix A.7), thus K  (x∗ )(x − x∗ )Y ≤ K(x) − K(x∗ )Y + γ2 x∗ − x2X and thus from (4.6) 2 Re(x∗ − x ˆ, x∗ − x)X ≤ 2wX K(x∗ ) − K(x))Y + wX γ x∗ − x2X for all x ∈ B(ˆ x, ρ) ∩ D(K) which proves the estimate (4.5) with σ = γwX < 1 and ϕ(t) ˜ = 2wX t for t ≥ 0. (c) Let Assumption 4.12 hold for ϕ(t) = βt for all t ≥ 0. For any fixed z ∈ X and t ∈ K sufficiently small such that x := x∗ − tz ∈ B(ˆ x, ρ) ∩ D(K) we substitute x into (4.5) for ϕ(t) = βt which yields   ≤ |t|2 σz2X + β K(x∗ − tz) − K(x∗ )Y 2 Re t (x∗ − x ˆ, z)X βγ 2 |t| z2X . ≤ |t|2 σz2X + |t|βK  (x∗ )zY + 2

128

Nonlinear Inverse Problems

  Dividing by |t| and letting t tend to zero yields2 (x∗ − x ˆ, z)X  ≤ β2 AzY for all ˆ belongs to the range of A∗ . z ∈ X where again A = K  (x∗ ). We show that x∗ − x ∗ ˆ is orthogonal to the nullspace N (A) of A. We choose First we note that x − x a singular system {μj , xj , yj : j ∈ J} for A, see Appendix A.6, Theorem A.57,

ˆ in the form x∗ − x ˆ = j∈J γj xj . where J is finite or J = N and expand x∗ − x ˆ is orthogonal (The component in the nullspace N (A) vanishes because x∗ − x to N (A).) We set Jn = J if J is finite and Jn = {1, . . . , n} if J = N. For γ ˆ, z)X  ≤ β2 Az is equivalent to z = j∈Jn μj2 xj the inequality (x∗ − x j  2 2 γj β 2 γj2 ≤ ; 2 μj 4 μ2j j∈Jn

j∈Jn

2 γ2 that is, j∈Jn μj2 ≤ β4 . This proves that w := j∈J j ˆ.  Finally, A∗ w = x∗ − x

γj μj

yj ∈ Y is well-defined.

Under the Assumptions 4.9 and 4.12, we are able to prove convergence and also rates of convergence as in the linear theory. As we know from the linear theory there are (at least) two √ strategies to choose the regularization parameter α. To achieve the rate O( δ), we should choose α = α(δ) to be proportional to δ (a priori choice) or such that the “discrepancy” K(xα(δ),δ ) − y δ Y to be proportional to δ (a posteriori choice). For nonlinear operators, K essentially the same arguments as for linear operators (substitute x = x ˆ and x = x∗ into (4.4)) show that lim sup K(xα,δ ) − y δ Y ≤ δ α→0

and

lim K(xα,δ ) − y δ Y = K(ˆ x) − y δ Y

α→∞

for any choice of minimizers xα,δ . However, for nonlinear operators, the mapping α → K(xα,δ ) − y δ Y is not necessarily continuous (see, e.g, [219] for a discussion of this topic). Therefore, the discrepancy principle is not well-defined unless more restrictive assumptions are made. In the following, we just take the possibility to choose the regularization parameter by the discrepancy principle as an assumption (see also [245], Section 4.1.2). ˆ. Theorem 4.14 Let Assumptions 4.9 and 4.12 hold and let ρ > 2x∗ − x (a) Let α = α(δ) be chosen such that c−

δ2 δ2 ≤ α(δ) ≤ c+ ϕ(δ) ϕ(δ)

for all δ > 0

where c+ ≥ c− > 0 are independent of δ (a priori choice), (b) or assume that there exists r+ > r− ≥ 1 and α(δ) > 0 such that r− δ ≤ K(xα(δ),δ ) − y δ Y ≤ r+ δ

for all δ > 0

(a posteriori choice) where xα(δ),δ denotes a minimum of the Tikhonov functional Jα(δ),δ on D(K). 2 Note

that the phases of t/|t| can be chosen arbitrarily!

4.2

The Nonlinear Tikhonov Regularization

129

Then xα(δ),δ ∈ B(ˆ x, ρ) for sufficiently small δ and   α(δ),δ ∗ 2 − x X = O ϕ(δ) and K(xα(δ),δ ) − y ∗ Y = O(δ) , x

δ → 0.

x, ρ) for sufficiently small δ. From Proof: We show first that xα(δ),δ ∈ B(ˆ (4.4) for x = x∗ , we conclude that ˆ2X K(xα,δ ) − y δ 2Y + αxα,δ − x

≤ ≤

δ2 ϕ(δ) ϕ(δ) ρ2 c− + 4

If c−

≤ α(δ) ≤ c+

δ2 ϕ(δ)

δ 2 + αx∗ − x ˆ2X ρ2 . δ2 + α 4

we conclude that xα(δ),δ − x ˆ2X ≤

δ2 α(δ)

(4.7a) (4.7b) +

ρ2 4

≤

2

≤ ρ for sufficiently small δ. If the discrepancy principle holds, then from (4.7b), 2 2 δ + α(δ)xα(δ),δ − x ˆ2X ≤ δ 2 + α(δ) r−

ρ2 , 4

2

and thus α(δ)xα(δ),δ − x ˆ2X ≤ α(δ) ρ4 because r− ≥ 1. This shows xα(δ),δ − x ˆX ≤ ρ and ends the first part of the proof. To show the error estimates, we use the variational source condition of Assumption 4.12 and (4.7a). K(xα,δ ) − y δ 2Y + αβxα,δ − x∗ 2X ≤ K(xα,δ ) − y δ 2Y   + αxα,δ − x ˆ2X − αx∗ − x ˆ2X + α ϕ K(x∗ ) − K(xα,δ )Y   ≤ δ 2 + α ϕ K(x∗ ) − K(xα,δ )Y   (4.8) ≤ δ 2 + α ϕ y δ − K(xα,δ )Y + δ . Let first α(δ) be chosen according to the discrepancy principle. Then   2 − 1) δ 2 + α(δ) β xα(δ),δ − x∗ 2X ≤ α(δ) ϕ (r+ + 1)δ (r− ≤

α(δ) (1 + r+ ) ϕ(δ)

where we have used that ϕ(sδ) ≤ sϕ(δ) for all s ≥ 1 (see Lemma A.73). This proves the assertion for xα(δ),δ − x∗ X after division by α(δ) and dropping the first term on the left hand side. The estimate for K(xα(δ),δ ) − y ∗ Y follows obviously from the triangle inequality because K(xα(δ),δ ) − y ∗ Y ≤ K(xα(δ),δ ) − y δ Y + δ ≤ (r+ + 1) δ. 2

2

δ δ Let now c− ϕ(δ) ≤ α(δ) ≤ c+ ϕ(δ) . Substituting this into (4.8) and dropping the second term on the left hand side yields

K(xα(δ),δ ) − y δ 2Y ≤ δ 2 +

 c+ δ 2  δ ϕ y − K(xα(δ),δ )Y + δ . ϕ(δ)

Now we set t = K(xα(δ),δ ) − y δ Y /δ for abbreviation. Then the previous estimate reads as  c+  ϕ (1 + t)δ ≤ 1 + c+ (1 + t) = 1 + c+ + c+ t , t2 ≤ 1 + ϕ(δ)

130

Nonlinear Inverse Problems

  where we used again ϕ (1 + t)δ ≤ (1 + t)ϕ(δ). Completing the square yields  c2 t ≤ c2+ + 1 + c+ + 4+ ; that is, K(xα(δ),δ ) − y δ Y ≤ cδ for some c > 0. Now we substitute this and the bounds of α(δ) into (4.8) again and drop the first term on the left hand side which yields  δ2 δ2  β xα(δ),δ − x∗ 2X ≤ δ 2 + c+ ϕ (1 + c)δ ≤ δ 2 + c+ (1 + c) δ 2 c− ϕ(δ) ϕ(δ) which yields c− β xα(δ),δ − x∗ 2X ≤ [1 + c+ (1 + c)] ϕ(δ) and ends the proof of the theorem.  By Theorem 4.13, the special case ϕ(t) = t corresponds to the √ source condition ˆ ∈ R (A∗ A)1/2 = R(A∗ ) and leads to the order O( δ) just as in the x∗ − x linear case. The cases ϕ(t) = tσ for σ > 1 are not covered by the previous theorem because these index functions are not concave anymore. For proving the analogue of Theorem 2.12 to get the optimal order of convergence up to O(δ 2/3 ), we have to use the classical source condition of Assumption 4.11, which is the obvious extension of the one in the linear case, see Theorem 2.12. A variational source condition for this case is not available. We follow the approach in [92] (see also [260] for the original proof). xX and (ii) Theorem 4.15 Let Assumptions 4.9 and 4.11 hold with ρ > 2x∗ −ˆ   ∗ ∗  ∗ σ/2   ∗ ∗  ∗ σ/2  ∗ modified in the way that x −ˆ x= K (x ) K (x ) v ∈ R K (x ) K (x )   ∗ ∗  ∗ −1/2 ∗ for some σ ∈ [1, 2] and v ∈ X such that γ K (x ) K (x ) (x − x ˆ) X < 1 where γ denotes the Lipschitz constant from part (i) of Assumption 4.11. We choose α(δ) such that c− δ 2/(σ+1) ≤ α(δ) ≤ c+ δ 2/(σ+1) . Then xα(δ),δ − x∗ X α(δ),δ

K(x

∗

) − y Y

  = O δ σ/(σ+1) ,   = O δ σ/(σ+1) ,

δ → 0, δ → 0.

Proof: We leave α > 0 arbitrary til the end of the proof. We set A := K  (x∗ ) and choose a singular system {μj , xj , yj : j ∈ J} for A, see Appendix A.6, ˆ as Theorem A.57, where J is finite or J = N. Then we write x∗ − x



x∗ − x ˆ = (A∗ A)σ/2 v = μσj vj xj = A∗ w with w = μσ−1 vj y j j j∈J

j∈J

where vj = (v, xj )X are the expansion coefficients of v. Then γwY < 1. We define zα ∈ X by (4.9) zα = x∗ − α (A∗ A + αI)−1 A∗ w ; that is, ∗

zα = x

− α

j∈J

μσj vj xj , thus zα − x∗ 2X = α2 2 μj + α j∈J



μσj μ2j + α

2 |vj |2 .

4.2

The Nonlinear Tikhonov Regularization

131

Later, we will also need the form 

μσ−1

 μσ+1 j j σ−1 ∗ A(zα − x ) + αw = α vj y j μj − 2 vj y j = α 2 μj + α μ2j + α j∈J

j∈J

with ∗

A(zα − x ) +

αw2Y

= α

4





j∈J

μσ−1 j μ2j + α

2 |vj |2 .

With the elementary estimate (see Problem 4.3) μt ≤ ct αt/2−1 , μ2 + α

μ ≥ 0,

for t = σ and t = σ − 1, respectively, (where ct depends on t only) we obtain zα − x∗ 2X A(zα − x∗ ) + αw2Y

≤ ≤

c2σ ασ , c2σ−1 ασ+1 .

(4.10a) (4.10b)

In particular, zα converges to x∗ as α → 0 and is, therefore, in B(ˆ x, ρ) ∩ D(K) for sufficiently small α. We use the optimality of xα,δ ; that is (4.4), for x = zα to obtain K(xα,δ ) − y δ 2Y + αxα,δ − x ˆ2X ≤ K(zα ) − y δ 2Y + αzα − x ˆ2X . With xα,δ − x ˆ2X zα − x ˆ2X

  = x∗ − x ˆ2X + 2 Re xα,δ − x∗ , x∗ − x ˆ X + xα,δ − x∗ 2X   = x∗ − x ˆ2X + 2 Re A(xα,δ − x∗ ), w Y + xα,δ − x∗ 2X ,   = x∗ − x ˆ2X + 2 Re A(zα − x∗ ), w Y + zα − x∗ 2X

we obtain   K(xα,δ ) − y δ 2Y + 2α Re w, A(xα,δ − x∗ ) Y + αxα,δ − x∗ 2X   ≤ K(zα ) − y δ 2Y + 2α Re w, A(zα − x∗ ) Y + αzα − x∗ 2X ,

and thus K(xα,δ ) − y δ + αw2Y + αxα,δ − x∗ 2X   ≤ α2 w2Y + 2α Re w, K(xα,δ ) − y δ − A(xα,δ − x∗ ) Y   +K(zα ) − y δ 2Y + 2α Re w, A(zα − x∗ ) Y + αzα − x∗ 2X . Now we use K(xα,δ ) = K(x∗) + A(xα,δ −x∗) + rα,δ = y ∗ + A(xα,δ −x∗ ) + rα,δ

and

K(zα ) = K(x∗ ) + A(zα − x∗ ) + sα = y ∗ + A(zα − x∗ ) + sα

132

Nonlinear Inverse Problems

with rα,δ Y ≤ γ2 xα,δ − x∗ 2X and sα Y ≤ γ2 zα − x∗ 2X and obtain K(xα,δ ) − y δ + αw2Y + αxα,δ − x∗ 2X     ≤ α2 w2Y + 2α Re w, y ∗ − y δ Y + 2α Re w, rα,δ Y   + y ∗ − y δ + A(zα − x∗ ) + sα 2Y + 2α Re w, A(zα − x∗ ) Y

+ αzα − x∗ 2X     = α2 w2Y + 2α Re w, y ∗ − y δ Y + 2α Re w, rα,δ Y   + y ∗ − y δ 2Y + 2 Re A(zα − x∗ ) + sα , y ∗ − y δ Y   +A(zα − x∗ ) + sα 2Y + 2α Re w, A(zα − x∗ ) Y + αzα − x∗ 2X   = 2 Re A(zα − x∗ ) + αw, y ∗ − y δ Y + 2 Re(sα , y ∗ − y δ )Y   + α2 w2Y + 2α Re w, rα,δ Y + A(zα − x∗ ) + sα 2Y   + 2α Re w, A(zα − x∗ ) + sα Y − 2α Re(w, sα )Y + αzα − x∗ 2X   = 2 Re A(zα − x∗ ) + αw, y ∗ − y δ Y + 2 Re(sα , y ∗ − y δ )Y   + 2α Re w, rα,δ Y + A(zα − x∗ ) + αw + sα 2Y − 2α Re(w, sα )Y + αzα − x∗ 2X

≤

2δA(zα − x∗ ) + αwY + γδzα − x∗ 2X + α γwY xα,δ − x∗ 2X + 2A(zα − x∗ ) + αw2Y γ2 zα − x∗ 4X + αγwY zα − x∗ 2X + αzα − x∗ 2X . + 2

Now we use that γwY < 1 and thus   K(xα,δ ) − y δ + αw2Y + α 1 − γwY xα,δ − x∗ 2X ≤

2δA(zα − x∗ ) + αwY + 2A(zα − x∗ ) + αw2Y γ2 zα − x∗ 4X + (γδ + αγwY + α)zα − x∗ 2X . + 2

So far, we have not used the definition of zα . We substitute the estimates (4.10a), (4.10b) and arrive at   K(xα,δ ) − y δ + αw2Y + α 1 − γwY xα,δ − x∗ 2X   ≤ c δ α(σ+1)/2 + ασ+1 + α2σ + δ ασ for some c > 0. Dropping one of the terms on the left hand side yields K(xα,δ ) − y δ 2Y

≤ ≤

xα,δ − x∗ 2X

≤

2K(xα,δ ) − y δ + αw2Y + 2α2 w2Y   c δ α(σ+1)/2 + ασ+1 + α2σ + δ ασ + α2 ,   c δ α(σ−1)/2 + ασ + α2σ−1 + δ ασ−1

4.2

The Nonlinear Tikhonov Regularization

133

for some c > 0. The choice c− δ 2/(σ+1) ≤ α(δ) ≤ c+ δ 2/(σ+1) yields the desired result (note that 1 ≤ σ ≤ 2).  We note that the modified defect K(xα(δ),δ )−y ∗ +α(δ)wY satisfies K(xα(δ),δ )− y ∗ + α(δ)wY ≤ cδ (see Problem 4.4). In the next subsection, we apply the result to the Tikhonov regularization of a parameter identification problem.

4.2.3

A Parameter-Identification Problem

Let f ∈ L2 (0, 1) be given. It is the aim to determine the parameter function c ∈ L2 (0, 1), c ≥ 0 on (0, 1), in the boundary value problem −u (t) + c(t) u(t) = f (t) , 0 < t < 1 ,

u(0) = u(1) = 0 ,

(4.11)

from perturbed data uδ (t). We recall the Sobolev spaces H p (0, 1) from (1.24) as the spaces p



H (0, 1) =

u∈C

p−1

[0, 1] : u

(p−1)

t (t) = α+

 ψ(s) ds , α ∈ R , ψ ∈ L (0, 1) 2

0

and set u(p) := ψ for the pth derivative. Note that H p (0, 1) ⊂ C[0, 1] for p ≥ 1 by definition. Then the differential equation of (4.11) for u ∈ H 2 (0, 1) is understood in the L2 −sense. It is an easy exercise (see Problem 4.7) to show that u∞ ≤ u L2 (0,1) for all u ∈ H 1 (0, 1) with u(0) = 0. First, we consider the direct problem and show that the boundary value problem is equivalent to an integral equation of the second kind. Lemma 4.16 Let f ∈ L2 (0, 1) and c ∈ L2+ (0, 1) := {c ∈ L2 (0, 1) : c ≥ 0 almost everywhere on (0, 1)}. (a) If u ∈ H 2 (0, 1) solves (4.11) then u solves the integral equation 1 u(t) +

1 g(t, s) c(s) u(s) ds =

0



where g(t, s) =

g(t, s) f (s) ds ,

t ∈ [0, 1] ,

(4.12)

0

s(1 − t), 0 ≤ s ≤ t ≤ 1, t(1 − s), 0 ≤ t ≤ s ≤ 1.

(b) We define the integral operator G : L2 (0, 1) → L2 (0, 1) by 

1 g(t, s) v(s) ds = (1 − t)

(Gv)(t) = 0

2

1

t

(1 − s) v(s) ds ,

s v(s) ds + t 0

t

t ∈ (0, 1), v ∈ L (0, 1). The operator G is bounded from L2 (0, 1) into H 2 (0, 1) and (Gv) = −v.

134

Nonlinear Inverse Problems

(c) If u ∈ C[0, 1] is a solution of (4.12); that is, of u + G(cu) = Gf , then u ∈ H 2 (0, 1) and u is a solution of (4.11). Note that the right-hand side of (4.12) is continuous because Gf ∈ H 2 (0, 1). Proof: (a) Let u ∈ H 2 (0, 1) solve (4.11) and set h = f − cu. Then h ∈ 2 L (0, 1) (because u is continuous) and −u = h. Integrating this equation twice and using the boundary conditions u(0) = u(1) = 0 yields the assertion (see Problem 4.5). (b) Let v ∈ L2 (0, 1) and set t u(t) = (Gv)(t) = (1 − t)

1

ψ(t) = −

1 (1 − s) v(s) ds ,

v(s) ds + 0

t ∈ [0, 1] ,

t

0

t

(1 − s) v(s) ds ,

s v(s) ds + t

t ∈ [0, 1] .

0

t Then it is easy to see that u(t) = 0 ψ(s) ds. Therefore, u ∈ H 1 (0, 1) and u = ψ. From the definitions of ψ and H 2 (0, 1), we observe that u ∈ H 2 (0, 1) and v = −u . (c) Let now u ∈ C[0, 1] be a solution of (4.12) and set again h = f − cu. Then again h ∈ L2 (0, 1), and u has the representation u = Gh. By part (b), we  conclude that u ∈ H 2 (0, 1) and −u = h = f − cu. Theorem 4.17 The integral equation (4.12) and the boundary value problem (4.11) are uniquely solvable for all f ∈ L2 (0, 1) and c ∈ L2+ (0, 1). Furthermore,  2 there exists  γ > 0 (independent of f and c) such that uH (0,1) ≤ γ 1 + cL2 (0,1) f L2 (0,1) . Proof:

By the previous lemma we have to study the integral equation u + G(cu) = Gf

(4.13)

with the integral operator G with kernel g from the previous lemma. The operator T : u → G(cu) is bounded from C[0, 1] into H 2 (0, 1) and thus compact from C[0, 1] into itself (see again Problem 4.5). Now we use the following result from linear functional analysis (see Theorem A.36 of the Appendix A.3): If the homogeneous linear equation u + T u = 0 with the compact operator T from C[0, 1] into itself admits only the trivial solution u = 0 then the inhomogeneous equation u + T u = h is uniquely solvable for all h ∈ C[0, 1], and the solution depends continuously on h. In other words, if I + T is one-to-one then also onto and I + T is boundedly invertible. Therefore, we have to show injectivity of I + T in C[0, 1]. Let u ∈ C[0, 1] solve (4.13) for h = 0. Then u ∈ H 2 (0, 1) and u solves (4.11) for f = 0 by the previous lemma. Multiplication of (4.11) by u(t) and integration yields 1 0 = 0

   −u (t) + c(t) u(t) u(t) dt =

1 0

  2  u (t) + c(t) u(t)2 dt ,

4.2

The Nonlinear Tikhonov Regularization

135

where we used partial integration and the fact that u vanishes at the boundary of [0, 1]. Since c ≥ 0 we conclude that u vanishes on [0, 1]. Therefore, u is constant and thus zero because of the boundary conditions. Therefore, I + T is one-to-one and thus invertible. This shows that (4.11) is uniquely solvable in H 2 (0, 1) for every f, c ∈ L2 (0, 1) with c ≥ 0 almost everywhere on (0, 1). In order to derive the explicit estimate for uH 2 (0,1) we observe first that any solution u ∈ H 2 (0, 1) of (4.11) satisfies u 2L2 (0,1) ≤ f L2 (0,1) uL2 (0,1) . Indeed, this follows by multiplication of the differential equation by u(t) and integration: 1 −

1



u (t) u(t) dt + 0

2

1 f (t) u(t) dt ≤ f L2 (0,1) uL2 (0,1) .

c(t) u(t) dt = 0

0

Partial integration and the assumption c(t) ≥ 0 yields the estimate u 2L2 (0,1) ≤ f L2 (0,1) uL2 (0,1) . With u∞ ≤ u L2 (0,1) and uL2 (0,1) ≤ u∞ this implies that u∞ ≤ f L2 (0,1) . Therefore uH 2 (0,1)

= G(f − cu)H 2 (0,1) ≤ GL(L2 (0,1),H 2 (0,1)) f − cuL2 (0,1)   ≤ GL(L2 (0,1),H 2 (0,1)) f L2 (0,1) + cL2 (0,1) u∞   ≤ 1 + cL2 (0,1) GL(L2 (0,1),H 2 (0,1)) f L2 (0,1) .

 We can even show existence and uniqueness for c from a small open neighbor hood of L2+ (0, 1) = c ∈ L2 (0, 1) : c ≥ 0 on (0, 1) . Corollary 4.18 There exists δ > 0 such that the boundary value problem (4.11) is uniquely solvable for all f ∈ L2 (0, 1) and c ∈ Uδ where   c1 ∈ L2+ (0, 1), h ∈ L2 (0, 1) 2  Uδ := c = c1 + h ∈ L (0, 1) : . 1 + c1 L2 (0,1) hL2 (0,1) < δ   Furthermore, there exists γ > 0 such that uH 2 (0,1) ≤ γ 1+cL2 (0,1) f L2 (0,1) for all f ∈ L2 (0, 1) and c ∈ Uδ . Proof: Let Kc1 : L2 (0, 1) → H 2 (0, 1) be the operator f → u where u is the solution of (4.11) for c1 ∈ L2+ (0, 1). Let c = c1 + h ∈ Uδ . We consider the fixed u) = Kc1 f for u ˜ ∈ H 2 (0, 1). We have for v ∈ H 2 (0, 1) point equation u ˜ + Kc1 (h˜ that Kc1 (hv)H 2 (0,1)

≤ ≤

γ(1 + c1 L2 (0,1) ) hvL2 (0,1) γ(1 + c1 L2 (0,1) ) hL2 (0,1) v∞ ≤ γ δvH 2 (0,1) .

For δ < 1/γ we observe that v → Kc1 (hv) is a contraction and, by the Conu) = Kc1 f has a unique solution u ˜ ∈ H 2 (0, 1) traction Theorem A.31, u ˜ + Kc1 (h˜

136

Nonlinear Inverse Problems

and ˜ uH 2 (0,1)

≤ ≤ ≤

1 + c1 L2 (0,1) 1 Kc1 f H 2 (0,1) ≤ γ f L2 (0,1) 1 − δγ 1 − δγ 1 + c1 + hL2 (0,1) + δ f L2 (0,1) γ 1 − δγ 1+δ (1 + cL2 (0,1) ) f L2 (0,1) γ 1 − δγ

because hL2 (0,1) ≤ δ. Finally, we note that the equation u ˜ = Kc1 f − Kc1 (h˜ u) u = f.  is equivalent to −˜ u + (c1 + h)˜ We note that the set Uδ is an open set containing L2+ (0, 1) (see Problem 4.5). Therefore, K can be extended to the set Uδ . As a next step towards the inverse problem, we show that the nonlinear mapping K : c → u is continuous and even differentiable. Theorem 4.19 Let Uδ ⊃ L2+ (0, 1) be as in the previous corollary and let K : Uδ → H 2 (0, 1) defined by K(c) = u where u ∈ H 2 (0, 1) solves the boundary value problem (4.11). Then K is continuous and even differentiable in every c ∈ Uδ . The derivative is given by K  (c)h = v where v ∈ H 2 (0, 1) is the solution of the boundary value problem −v  (t) + c(t) v(t) = −h(t) u(t) , 0 < t < 1 ,

v(0) = v(1) = 0 ,

(4.14)

and u ∈ H 2 (0, 1) is the solution of (4.11) for c; that is, u = K(c). ˜ = K(c + h). Then u Proof: Let h ∈ L2 (0, 1) such that c + h ∈ Uδ and let u and u ˜ satisfy u = f, −˜ u +(c+h)˜

−u +(c+h)u = f +hu ,

u ˜(0) = u ˜(1) = u(0) = u(1) = 0 ,

u−u) = respectively. We subtract both equations which yields −(˜ u−u) +(c+h)(˜ −hu. The stability estimate yields   ˜ u − uH 2 (0,1) ≤ γ 1 + c + hL2 (0,1) huL2 (0,1)   ≤ γ 1 + cL2 (0,1) + δ hL2 (0,1) u∞ which proves continuity (even Lipschitz continuity on bounded sets for c). For the differentiability, we just subtract the equations for u and v from the one for u ˜ u − u − v) = −hv , (4.15) −(˜ u − u − v) + (c + h)(˜ with homogeneous boundary conditions. The stability estimate yields   ˜ u − u − vH 2 (0,1) ≤ γ 1 + c + hL2 (0,1) hvL2 (0,1)   ≤ γ 1 + cL2 (0,1) + δ hL2 (0,1) v∞ .

4.2

The Nonlinear Tikhonov Regularization

137

The stability estimate applied to v yields   v∞ ≤ vH 2 (0,1) ≤ γ 1 + cL2 (0,1) huL2 (0,1)   ≤ γ 1 + cL2 (0,1) hL2 (0,1) u∞ which altogether ends up to  2 ˜ u − u − vH 2 (0,1) ≤ γ 2 1 + cL2 (0,1) + δ h2L2 (0,1) u∞ . 

This proves differentiability.

In the following, we consider the parameter-to-solution map K as a mapping from Uδ ⊂ L2 (0, 1) into L2 (0, 1) instead of H 2 (0, 1). Then, of course, K is also differentiable with respect to this space with the same derivative. From the theory, we know that for Assumption 4.11, we need the adjoint of K  (c). Lemma 4.20 The adjoint operator K  (c)∗ : L2 (0, 1) → L2 (0, 1) is given by K  (c)∗ w = −u y where u = K(c), and y ∈ H 2 (0, 1) solves the following boundary value problem (the “adjoint problem”): −y  (t) + c(t) y(t) = w(t) , 0 < t < 1 ,

y(0) = y(1) = 0 ,

(4.16)

Proof: Let h, w ∈ L2 (0, 1) and v the solution of (4.14) for c = c∗ and y the solution of (4.16). By partial integration we compute    K (c)h, w L2 (0,1)

1 =

v(t) [−y  (t) + c(t) y(t)] dt

0

1



1

[−v (t) + c(t) v(t)] y(t) dt = −

= 0

h(t) u(t) y(t) dt 0

= −(h, uy)L2 (0,1) . This proves the assertion.



Now we can formulate condition (ii) of Assumption 4.11 for a linear index function ϕ: The existence of w ∈ L2 (0, 1) with c∗ − cˆ = K  (c∗ )∗ w is equivalent to the existence of y ∈ H 2 (0, 1) with y(0) = y(1) = 0 and c∗ − cˆ = −u∗ y. Therefore, the condition is equivalent to c∗ − cˆ ∈ H 2 (0, 1) ∩ H01 (0, 1) . u∗ This includes smoothness of c∗ − cˆ as well as a sufficiently strong boundary condition because also u∗ vanishes at the boundary of [0, 1].

138

Nonlinear Inverse Problems

4.2.4

A Glimpse on Extensions to Banach Spaces

Recalling the classical Tikhonov functional Jα,δ in Hilbert spaces from (4.2) we observe that the first term measures the misfit in the equation while the second part serves as a penalty term. The error xα,δ − x∗ in the solution is measured in a third metric. In many cases, the canonical space for the unknown quantity x is only a Banach space rather than a Hilbert space. For example, in parameter identification problems as in the previous subsection the canonical space for the parameters are L∞ −spaces rather than L2 −spaces. The Hilbert space setting in the previous subsection only works because the pointwise multiplication is continuous as a mapping from L2 × H 2 into L2 . For more general partial differential equation this is not always true. For an elaborate motivation why to use Banach space settings, we refer to Chapter I of the excellent monograph [245] by Schuster, Kaltenbacher, Hofmann, and Kazimierski. In this subsection, we will get the flavor of some aspects of this theory. Let X and Y be Banach spaces, K : X ⊃ D(K) → Y a (nonlinear) operator with the domain of definition D(K) ⊂ X where D(K) is again convex and closed, and let x∗ ∈ D(K) be the exact solution of K(x) = y ∗ for some y ∗ . In the following we fix this pair x∗ , y ∗ . As before, y ∗ is perturbed by y δ ∈ Y such that y δ − y ∗ Y ≤ δ for all δ ∈ (0, δmax ) for some δmax > 0. We note that we measure the error in the data with respect to the Banach space norm. The penalty term x − x ˆ2X is now replaced by any convex and continuous function Ω : X → [0, ∞] where D(Ω) := {x ∈ X : Ω(x) < ∞} is not empty and, even more, D(K) ∩ D(Ω) = ∅. Further assumptions on Ω and K are needed to ensure the existence of minima of the Tikhonov functional Jα,δ (x) := K(x) − y δ pY + α Ω(x) ,

x ∈ D(K) ∩ D(Ω) .

(4.17)

Here p > 1 is a fixed parameter. Instead of posing assumptions concerning the weak topology, we just make the same assumption as at the beginning of Subsection 4.2.2. Assumption 4.21 The Tikhonov functional Jα,δ possesses global minima xα,δ ∈ D(K) ∩ D(Ω) on D(K) ∩ D(Ω) for all α, δ > 0; that is, K(xα,δ ) − y δ pY + α Ω(xα,δ ) ≤ K(x) − y δ pY + α Ω(x)

(4.18)

for all x ∈ D(K) ∩ D(Ω). It remains to specify the metric in which we measure the error in x. As we will see in a moment the norm in X is not always the best possibility. To have more flexibility, we take any “measure function” E(x) which measures the distance of x to x∗ . We only require that E(x) ≥ 0 for all x ∈ X and E(x∗ ) = 0. Then the variational source condition of Assumption 4.12 is generalized into the following form.

4.2

The Nonlinear Tikhonov Regularization

139

Assumption 4.22 (Variational Source condition) Let x∗ ∈ D(K) ∩ D(Ω) be a solution of K(x) = y ∗ . Furthermore, let δmax = sup K(x∗ ) − K(x)Y : x ∈ D(K) ∩ D(Ω) and ϕ : [0, δmax ) → R be a concave index function (see Definition 4.10), and, for some constants β > 0 and ρ > 0 let the following estimate hold.   β E(x) ≤ Ω(x) − Ω(x∗ ) + ϕ K(x∗ ) − K(x)Y  for all x ∈ Mρ := x ∈ D(K) ∩ D(Ω) : Ω(x) ≤ Ω(x∗ ) + ρ . Theorem 4.23 Let Assumptions 4.21 and 4.22 hold. (a) Let α = α(δ) be chosen such that c−

δp δp ≤ α(δ) ≤ c+ ϕ(δ) ϕ(δ)

for all δ ∈ (0, δ1 )

(4.19)

where c+ ≥ c− ≥ ρ1 are independent of δ and where δ1 ∈ (0, δmax ) is chosen such that ϕ(δ1 ) ≤ ρ c− (a priori choice), (b) or assume that there exists r+ > r− ≥ 1 and α(δ) > 0 such that r− δ ≤ K(xα(δ),δ ) − y δ Y ≤ r+ δ

for all δ ∈ (0, δ1 )

(4.20)

where δ1 ∈ (0, δmax ) is arbitrary (a posteriori choice). Then xα(δ),δ ∈ Mρ and   E(xα(δ),δ ) = O ϕ(δ)

and

K(xα(δ),δ ) − y ∗ Y = O(δ) ,

δ → 0.

Proof: We follow almost exactly the proof of Theorem 4.14. Substituting x = x∗ into (4.18) yields K(xα,δ ) − y δ pY + α Ω(xα,δ ) ≤ δ p + α Ω(x∗ ) .

(4.21)

First we show that xα(δ),δ ∈ Mρ . If α(δ) is chosen as in (4.19) then Ω(xα(δ),δ ) ≤

δp 1 1 + Ω(x∗ ) ≤ ϕ(δ) + Ω(x∗ ) ≤ ϕ(δ1 ) + Ω(x∗ ) α(δ) c− c−

which shows xα(δ),δ ∈ Mρ by the choice of δ1 . If α(δ) is chosen by the discrepancy principle (4.20) then again from (4.21) p − 1) δ p + α(δ) Ω(xα(δ),δ ) ≤ α(δ) Ω(x∗ ) (r−

and thus Ω(xα(δ),δ ) ≤ Ω(x∗ ) because r− ≥ 1. Therefore, xα(δ),δ ∈ M0 ⊂ Mρ . Now we show the rates of convergence. Applying the variational source condition we conclude from (4.21) that   K(xα,δ ) − y δ pY + αβ E(xα,δ ) ≤ δ p + α ϕ K(xα,δ ) − y ∗ Y

140

Nonlinear Inverse Problems   ≤ δ p + α ϕ K(xα,δ ) − y δ Y + δ .

(4.22)

If α = α(δ) is chosen according to the discrepancy principle (4.20) then   p (r− − 1)δ p + α(δ) β E(xα(δ),δ ) ≤ α(δ) ϕ (r+ + 1)δ   and thus β E(xα(δ),δ ) ≤ ϕ (r+ + 1)δ . Now we use the elementary estimate ϕ(sδ) ≤ s ϕ(δ) for all s ≥ 1 and δ ≥ 0 (see Lemma A.73 of Appendix A.8). Therefore, β E(xα(δ),δ ) ≤ (r+ + 1) ϕ(δ) . This proves the estimate for E(xα(δ),δ ). The estimate for K(xα(δ),δ ) − y ∗ Y follows obviously from the discrepancy inequality and the triangle inequality. Let now α = α(δ) be given by (4.19). From (4.22), we obtain, using the upper estimate of α(δ),  δp  ϕ K(xα,δ ) − y δ Y + δ . K(xα(δ),δ ) − y δ pY ≤ δ p + c+ ϕ(δ) We set t = K(xα(δ),δ ) − y δ Y /δ for abbreviation. Then the previous formula takes the form   ϕ δ(t + 1) ≤ 1 + c+ (t + 1) = (1 + c+ ) + c+ t tp ≤ 1 + c+ ϕ(δ) where we used the estimate ϕ(sδ) ≤ sϕ(δ) for s ≥ 1 again. Choose c > 0 with c (cp−1 − c+ ) > 1 + c+ . Then t ≤ c. Indeed, if t > c then tp − c+ t = t (tp−1 − c+ ) > c (cp−1 − c+ ) > 1 + c+ , a contradiction. This proves that K(xα(δ),δ ) − y δ Y ≤ c δ. Now we substitute this into the right-hand side of (4.22), which yields   1 δp + ϕ (c + 1)δ ≤ β E(xα(δ),δ ) ≤ ϕ(δ) + (c + 1) ϕ(δ) . α(δ) c− This ends the proof.



We note that in the case of Hilbert spaces X and Y and Ω(x) = x − x ˆ2X and ∗ 2 E(x) = x − x X Assumption 4.22 and Theorem 4.23 reduce to Assumption 4.12 and Theorem 4.14, respectively. Before we continue with the general theory, we apply this theorem to the special situation to determine a sparse approximation of the linear problem Kx = y δ . By sparse we mean that the solution x∗ ∈ X can be expressed by only finitely ˜ be a Banach space many elements of a given basis of X. Therefore, let X having a Schauder basis {bj : j ∈ N} with bj X˜ = 1 for all j ∈ N; that is, ˜ has a unique representation as x = every element x ∈ X j∈N xj bj where, of ˜ course, the convergence is understood in the norm of X. We define the subspace ˜ by X⊂X ⎧ ⎫ ⎨ ⎬



X = x= xj bj : |xj | < ∞ ⎩ ⎭ j∈N

j∈N

4.2

The Nonlinear Tikhonov Regularization

141



with the norm j∈N xj bj X := |xj | for x ∈ X. Then X is boundj∈N ˜ because x ˜ = ˜ ≤ edly imbedded in X ˜ = j∈N xj bj X j∈N |xj |bj X X

|x | = x for x ∈ X. Obviously, the space X is norm-isomorphic j X j∈N

∞ to the space 1 of sequences (xj ) such that j=1 |xj | converge, equipped with

∞ the canonical norm x 1 = j=1 |xj | for x = (xj ). As the Schauder basis of  (j)

1 we take e(j) : j = 1, 2, . . . ⊂ 1 where e(j) ∈ 1 is defined as ek = 0 for (j) k = j and ej = 1. Therefore, we can take directly 1 as the space X. We note that the dual of 1 is just ( 1 )∗ = ∞ , the space

of bounded sequences with the sup-norm and the dual pairing3 y, x ∞ , 1 = j∈N yj xj for y ∈ ∞ and x ∈ 1 . Also we note that 1 itself is the dual of the space c0 of sequences converging to zeros (see Example A.21). Therefore, by Theorem A.77 of Appendix A.9 the unit ball in 1 is weak∗ compact which is an important ingredient to prove existence of minimizers of the Tikhonov functional. With these introductory remarks, we are able to show the following result where we followed the presentation in [96]. Theorem 4.24 Let Y be a Banach space and K : 1 → Y be a linear bounded operator such that μ, Ke(j) Y ∗ ,Y → 0 as j → ∞ for all μ ∈ Y ∗ where μ, yY ∗ ,Y denotes the application of μ ∈ Y ∗ to y ∈ Y . Let Kx∗ = y ∗ and y δ ∈ Y with y δ − y ∗ Y ≤ δ. (a) Let xα,δ ∈ 1 be a minimizer of Jα,δ (x) = Kx − y δ pY + αx 1 ,

x ∈ 1 .

Then xα,δ ∈ 1 is sparse; that is, the number of non-vanishing components = 0 is finite. xα,δ j (b) For every j ∈ N let there exists fj ∈ Y ∗ with e(j) = K ∗ fj where K ∗ : Y ∗ → ( 1 )∗ = ∞ is the dual operator corresponding to K. Define the function ϕ : [0, ∞) → R as n  



. |x∗j | where γn = sup s f ϕ(t) := 2 inf γn t + j j n∈N

j>n

sj ∈{0,1,−1} j=1

Y∗

Then ϕ is a concave index function. With the choices (4.19) or (4.20) of α = α(δ) the following convergence rates hold:   xα(δ),δ − x∗  1 = O ϕ(δ) , Kxα(δ),δ − y ∗ Y = O(δ) as δ tends to zero. (c) If x∗ ∈ 1 is such that

∞

γjσ |x∗j | < ∞ for some σ > 0 then we have   xα(δ),δ − x∗  1 = O δ σ/(1+σ) , δ → 0 . j=1

3 Note that we denote the dual pairing by , x ∗ ∗ X ,X = (x) for  ∈ X and x ∈ X. The mapping (, x) → , xX ∗ ,X is bilinear.

142

Nonlinear Inverse Problems If x∗ is sparse or if (γn ) is bounded then we have xα(δ),δ − x∗  1 = O(δ) as δ tends to zero.

Proof: as

(a) Set z = Kxα,δ − y δ for abbreviation. The optimality of xα,δ reads

  z + KhpY − zpY ≥ −α xα,δ + h 1 − xα,δ  1 for all h ∈ 1 .

(4.23)

Define the sets A, B ⊂ R × Y as follows:  A = (r, y) ∈ R × Y : r > z + ypY − zpY ,    B = (r, Kh) ∈ R × Y : h ∈ 1 , r ≤ −α xα,δ + h 1 − xα,δ  1 . Then it is not difficult to show that A and B are convex, A is open, and A∩B = ∅ because of (4.23). Now we apply the separation theorem for convex sets (see Theorem A.69 of Appendix A.8). There exists (s, μ) ∈ R × Y ∗ and γ ∈ R such that (s, μ) = (0, 0) and s r + μ, yY ∗ ,Y ≥ γ ≥ s r + μ, KhY ∗ ,Y for all (r, y) ∈ A and (r , Kh) ∈ B . Letting r tend to infinity while keeping the other variables constant yields s ≥ 0. It is s = 0 because otherwise we would have μ, yY ∗ ,Y ≥ γ for all y ∈ Y (set r := z + ypY − zpY + 1) which would yield that also μ vanishes4 , a contradiction. Therefore, s > 0, and without loss of generality, s = 1. Now we set y = 0 and fix any h ∈  1 and let r tend to zero from above and set r = −α xα,δ + h 1 − xα,δ  1 . This yields the inequality   0 ≥ −α xα,δ + h 1 − xα,δ  1 + μ, KhY ∗ ,Y for all h ∈ 1 . For any t ∈ R and k ∈ N, we set h = te(k) and arrive at  α,δ       − t μ, Ke(k) Y ∗ ,Y ≥ 0 . α xα,δ k + t − xk   For fixed k with xα,δ = 0 we choose |t| so small such that sign xα,δ +t = k k   . Then the previous inequality reads as sign xα,δ k     t α sign xα,δ − μ, Ke(k) Y ∗ ,Y ≥ 0 . k   Choosing t > 0 and t < 0 yields α sign xα,δ = μ, Ke(k) Y ∗ ,Y ; that is, k    (k) μ, Ke Y ∗ ,Y  = α. This holds for every k ∈ J := k ∈ N : xα,δ = 0 . k This implies that J is finite because of μ, Ke(k) Y ∗ ,Y → 0 as k → ∞. (b) We apply Theorem 4.23 and have to verify Assumptions 4.21 and 4.22 for the special case D(K) = D(Ω) = X = 1 , E(x) = x − x∗  1 , and Ω(x) = x 1 for x ∈ X = 1 . First we show that ϕ is a concave index function. Indeed, ϕ is continuous, monotonic, and concave as the infimum of affine functions 4 the

reader should prove this himself.

4.2

The Nonlinear Tikhonov Regularization

143

(see Problem 4.6). Furthermore, ϕ(0) = 2 inf

n∈N

j>n

|x∗j | = 0 which shows

that ϕ is a concave index function. Assumption 4.21; that is, existence of a minimum of Jα,δ can be shown using results on weak- and weak∗-topologies. (The assumption that y ∗ , Ae(k) Y ∗ ,Y tends to zero for all y ∗ is equivalent to the weak∗-weak continuity of K. Then one uses that  · pY and  ·  1 are lower weak semi-continuous and lower weak∗ semi-continuous, respectively.) We do not carry out this part but refer to, e.g., [96]. Assumption 4.22 is therefore equivalent to (with β = 1)   for all x ∈ 1 . (4.24) x − x∗  1 − x 1 + x∗  1 ≤ ϕ K(x − x∗ )Y To prove this we have for any n ∈ N n

|xj − x∗j |

=

j=1

n

j=1

=

n

j=1

sj e(j) , x − x∗  ∞ , 1 =

n

sj K ∗ fj , x − x∗  ∞ , 1

j=1



n sj fj , K(x − x∗ )Y ∗ ,Y ≤ s f j j j=1

Y

∗

K(x − x∗ )Y

∗

≤ γn K(x − x )Y where sj = sign(xj − x∗j ) ∈ {0, 1, −1}. Therefore, with |xj | ≥ |x∗j | − |xj − x∗j |, n



|xj − x∗j | − |xj | + |x∗j |



≤ 2

n

|xj − x∗j | ≤ 2γn K(x − x∗ )Y .

j=1

j=1

Furthermore,



  |xj − x∗j | − |xj | + |x∗j | ≤ 2 |x∗j | ; j>n

that is,

j>n

⎡

x − x∗  1 − x 1 + x∗  1 ≤ 2 ⎣γn K(x − x∗ )Y +



⎤ |x∗j |⎦ (4.25)

j>n

This shows (4.24) since this estimate holds for all n ∈ N. Therefore, all of the assumptions of Theorem 4.23 are satisfied, and the error estimates are shown. (c) We estimate ϕ(t). First we observe that (γn ) is monotonically increasing. If γn is bounded then it converges to some finite γ ∈ R. Therefore, we let n tend to infinity in the definition of ϕ and arrive at ϕ(t) = 2γ t. We consider now the case that γn tends to infinity. First we estimate

1 σ ∗ c γn t + |x∗j | ≤ γn t + σ γj |xj | ≤ γn t + σ γn+1 j>n γn+1 j>n with c =

∞

j=1

γjσ |x∗j |. For sufficiently small t, the index   1 n(t) = max n ∈ N : γn ≤ 1/(1+σ) t

144

Nonlinear Inverse Problems

is well-defined and finite. Then γn(t) ≤ t−1/(1+σ) and γn(t)+1 > t−1/(1+σ) . Therefore, c ϕ(t) ≤ γn(t) t + σ ≤ (1 + c) tσ/(1+σ) . γn(t)+1 Finally, if x∗ is sparse exists n ∈ N such that x∗j = 0 for all j > n.

then there ∗ For that n the series j>n |xj | in (4.25) vanishes which shows the result for the  linear index function ϕ(t) = 2γn t. Remark: The reciprocals 1/γn play the role of the singular values in the case of Hilbert spaces X = 2 and Y with a singular system {μj , e(j) , g (j) : j ∈ J}. Indeed, the assumption e(j) = K ∗ fj is satisfied with fj = μ1j g (j) and for γn one has the form

2 n n



n s2j 1 γn2 = sup s f = sup = j j μ2 μ2 sj ∈{0,1,−1} j=1

The condition that   x∗ ∈ R (A∗ A)σ/2 .

∞

j=1

Y

sj ∈{0,1,−1} j=1

j

j=1

j

γjσ |x∗j | converges corresponds to the source condition

We go now back to the general case studied in Theorem 4.23. We wish to carry over Theorem 4.13 which proves the variational source condition from the classical one of Assumption 4.11. The essential formula used in the proof of part (a) of Theorem 4.13 was based on the binomial formula; that is, ˆ2X − x∗ − x ˆ2X − 2 (x − x∗ , x∗ − x ˆ)X . x∗ − x2X = x − x If we denote the penalty term by Ω(x); that is, Ω(x) = x − x ˆ2X then we can write this formula as x∗ − x2X = B Ω (x, x∗ ) := Ω(x) − Ω(x∗ ) − Ω (x∗ ), x − x∗ X ∗ ,X where Ω (x∗ ) : X → R is the Fr´echet derivative of Ω at x∗ and  , zX ∗ ,X is the application of ∈ X ∗ to z ∈ X. The function B Ω (x, x∗ ) is the famous Bregman distance corresponding to the function Ω. This can be extended to Banach spaces because for convex and differentiable functions Ω from a Banach space X into R the function B Ω (x, x∗ ) := Ω(x) − Ω(x∗ ) − Ω (x∗ ), x − x∗ X ∗ ,X ,

x∈X,

is nonnegative on X (see Lemma A.70 of Appendix A.8). If Ω is even strictly convex, then B Ω (x, x∗ ) = 0 ⇔ x = x∗ . If Ω : X → R is convex and only continuous then the subdifferential ∂Ω(x∗ ) ⊂ X ∗ ; that is, set of subgradients is non-empty (see Lemma A.72 of Appendix A.8). We recall that ∂Ω(x∗ ) ⊂ X ∗ is the set of all ∈ X ∗ with Ω(x) − Ω(x∗ ) −  , x − x∗ X ∗ ,X ≥ 0 for all x ∈ X .  If Ω is convex and differentiable then ∂Ω(x∗ ) = Ω (x∗ ) (Lemma A.72). Therefore, we formulate the following definition:

4.2

The Nonlinear Tikhonov Regularization

145

Definition 4.25 Let X be a normed space, A ⊂ X convex and open, and Ω : A → R convex and continuous with subdifferential ∂Ω(x∗ ) at some x∗ ∈ A. For

∈ ∂Ω(x∗ ) the Bregman distance is defined as B Ω (x, x∗ ) := Ω(x) − Ω(x∗ ) −  , x − x∗ X ∗ ,X ,

x ∈ A.

We note that the Bregman distance depends on the function Ω. It measures the defect of the function with its linearization at x∗ . With the Bregman distance as E(x) we have an analogue of Theorem 4.13. Lemma 4.26 Let x∗ ∈ D(K) ∩ D(Ω) be a solution of K(x) = y ∗ and let ϕ : [0, ∞) → R be a concave index function. Furthermore, let Ω : X → R be convex and continuous and ∈ ∂Ω(x∗ ) and E(x) = B Ω (x, x∗ ). (a) Then, for this particular choice of E(x), the variational source condition of Assumption 4.22 is equivalent to the following condition: There exists ρ > 0 and 0 ≤ σ < 1 such that    , x∗ − xX ∗ ,X ≤ σ B Ω (x, x∗ ) + ϕ K(x∗ ) − K(x)Y (4.26) for all x ∈ Mρ . (b) Assume that there exist w ∈ Y ∗ and a concave index function η such that

= K  (x∗ )∗ w and   K  (x∗ )(x − x∗ )Y ≤ η K(x∗ ) − K(x)Y for all x ∈ Mρ . Then Assumption 4.22 holds with ϕ(t) = wY ∗ η(t) for t ≥ 0. Proof: (a) This follows directly from the definitions of E(x) and B Ω (x, x∗ ). Indeed, the estimate in Assumption 4.22 reads as   β B Ω (x, x∗ ) ≤ Ω(x) − Ω(x∗ ) + ϕ K(x∗ ) − K(x)Y ; that is,   β B Ω (x, x∗ ) ≤ B Ω (x, x∗ ) +  , x − x∗ X ∗ ,X + ϕ K(x∗ ) − K(x)Y which is equivalent to    , x∗ − xX ∗ ,X ≤ (1 − β) B Ω (x, x∗ ) + ϕ K(x∗ ) − K(x)Y . This proves part (a) with σ = 1 − β if β ≤ 1 and σ = 0 otherwise. (b) If = K  (x∗ )∗ w then  , x∗ − xX ∗ ,X

= K  (x∗ )∗ w, x∗ − xX ∗ ,X = w, K  (x∗ )(x∗ − x)Y ∗ ,Y   ≤ wY ∗ K  (x∗)(x∗ − x)Y ≤ wY ∗ η K(x∗ ) − K(x)Y .

This proves the condition of (a) with σ = 0; that is, β = 1.



Combining Theorem 4.23 with these particular choices, we have the following theorem:

146

Nonlinear Inverse Problems

Theorem 4.27 Let Assumption 4.21 hold and let x∗∈ D(K) ∩ D(Ω) be a solution of K(x) = y ∗ . Furthermore, let δmax = sup K(x∗ ) − K(x)Y : x ∈ D(K) ∩ D(Ω) and ϕ : [0, δmax ) → R be a concave index function, and for some constants ρ > 0 and 0 ≤ σ < 1, let the source condition (4.26) hold for some

∈ ∂Ω(x∗ ). Let α = α(δ) be chosen according to (4.19) or (4.20). Then we have the error estimates   B Ω (xα(δ),δ , x∗ ) = O ϕ(δ) and K(xα(δ),δ ) − y ∗ Y = O(δ) as δ tends to zero. As a particular and obviously important example, we now take Ω(x) := x − x ˆpX ,

x∈X,

for some p > 1. Then one would like to characterize the Bregman distance — or, at least, construct lower bounds of B Ω (x, x∗ ) in terms of x − x∗ pX . This leads to the concept of p−convex Banach spaces. Definition 4.28 A Banach space X is called p−convex for some p > 0 if there exists c > 0 such that   x + ypX − xpX −  x , yX ∗ ,X ≥ c ypX for all x ∈ ∂  · pX (x) and x, y ∈ X.5 In other words, for p−convex spaces the Bregman distance B Ωz (x, z) corresponding to Ω(x) = x − x ˆpX can be bounded below by c x − zpX for all x, z ∈ X. Therefore, if the assumptions of   the previous Theorem 4.27 holds one has the rate xα(δ),δ − x∗ pX = O ϕ(δ) . As a particular example, we show that Lp (D) are p−convex for all p > 2. Lemma 4.29 Let p > 1 and D ⊂ Rn open and    x(t)p dt = f (x) = xp p L (D)

for x ∈ Lp (D) .

D

(a) Then f is differentiable and  p−1  sign x(t) dt f  (x)y = p y(t) x(t)

for x, y ∈ Lp (D) .

D

(b) Let p > 2. Then there exists cp > 0 with f (x + y) − f (x) − f  (x)y ≥ cp ypLp (D)

for all x, y ∈ Lp (D) .

5 Actually, the classical definition uses the dual mapping instead of the subdifferential. However, by Asplund’s theorem (see [48]) they are equivalent.

4.2

The Nonlinear Tikhonov Regularization

147

Proof: (a) First we observe that the integral for f  (x)y exists by H¨older’s p and s = p then 1r + 1s = 1 and thus inequality. Indeed, set r = p−1 

|x(t)|p−1 |y(t)| dt

D

≤

⎛ ⎞1/r ⎛ ⎞1/s   ⎝ |x(t)|r(p−1) dt⎠ ⎝ |y(t)|s dt⎠ D

D

⎛ ⎞(p−1)/p ⎛ ⎞1/p   ⎝ |y(t)|p dt⎠ = ⎝ |x(t)|p dt⎠ D

D

p = xp−1 Lp (D) yL (D) .

We use the following elementary estimate. There exist constants c+ > 0 and cp ≥ 0 with cp = 0 for p ≤ 2 and cp > 0 for p > 2 such that for all z ∈ R ' c+ |z|p if p ≤ 2 or |z| ≥ 12 , p p cp |z| ≤ |1 + z| − 1 − p z ≤ (4.27) c+ |z|2 if p > 2 and |z| ≤ 12 . We give a proof of this estimate in Lemma A.74 for the convenience of the reader. Therefore, for any x, y ∈ R with x = 0, we have  y p y  ≥ 0 |x + y|p − |x|p − p y |x|p−1 sign x = |x|p 1 +  − 1 − p x x and

 y p y  |x + y|p − |x|p − p y |x|p−1 sign x = |x|p 1 +  − 1 − p x x ' c+ |x|p |y/x|p = c+ |y|p if p ≤ 2 or 2|y| ≥ |x| , ≤ c+ |x|p |y/x|2 = c+ |x|p−2 |y|2 if p > 2 and 2|y| ≤ |x| .

This holds obviously also for x = 0. Now we apply this to x(t) and y(t) with x, y ∈ Lp (D). This shows already that         x(t) + y(t)p − x(t)p − p y(t) x(t)p−1 sign x(t) dt ≥ 0 . D

Next we show that         x(t) + y(t)p − x(t)p − p y(t) x(t)p−1 sign x(t) dt ≤ c ymin{2,p} Lp (D) D

for yLp (D) ≤ 1. This would finish the proof of part (a) because p > 1.  For proving this estimate we define T := t ∈ D : 2|y(t)| ≥ |x(t)| . Then         x(t) + y(t)p − x(t)p − p y(t) x(t)p−1 sign x(t) dt ≤ c+ |y(t)|p dt T

T

148

Nonlinear Inverse Problems

and



      x(t) + y(t)p − x(t)p − p y(t) x(t)p−1 sign x(t) dt

D\T

≤



⎧ ⎪ ⎪ ⎪ ⎪ ⎨

|y(t)|p dt

if p ≤ 2 ,

|x(t)|p−2 |y(t)|2 dt

if p > 2 .

c+

⎪ ⎪ c+ ⎪ ⎪ ⎩



D\T

D\T

If p ≤ 2 we just add the two estimates and have shown the estimate        x(t) + y(t)p − x(t)p − p y(t) x(t)p−1 sign x(t) dt ≤ c+ yp p

L (D)

D

If p > 2 we apply H¨ older’s inequality to the integral p Indeed, set r = p−2 and s = p2 then 1r + 1s = 1 and thus 

|x(t)|p−2 |y(t)|2 dt

 ≤

D\T

1/r

D\T

⎛ ⎜ = ⎝



|x(t)|r(p−2) dt

≤

⎜ ⎝

⎜ ⎝

D\T 2 xp−2 Lp (D) yLp (D)

D\T

⎛

⎞(p−2)/p ⎛ ⎟ |x(t)|p dt⎠



|x(t)|p−2 |y(t)|2 dt. ⎞1/s



⎟ |y(t)|2s dt⎠

D\T



.

⎞2/p

⎟ |y(t)|p dt⎠

D\T

.

Therefore, 

      x(t) + y(t)p − x(t)p − p y(t) x(t)p−1 sign x(t) dt

D

≤

2 c+ ypLp (D) + c+ xp−2 Lp (D) yLp (D) .

For p > 2 and yLp (D) ≤ 1 we have ypLp (D) ≤ y2Lp (D) and thus        x(t) + y(t)p − x(t)p − p y(t) x(t)p−1 sign x(t) dt ≤ c y2 p L (D) . D

This proves part (a). (b) By part (a) we have to show that            x(t) + y(t)p − x(t)p − p y(t) x(t)p−1 sign x(t) − cp y(t)p dt ≥ 0 D

4.3

The Nonlinear Landweber Iteration

149

for all x, y ∈ Lp (D). We show that the integrand is non-negative. Indeed, we have  y p   y p y    |x + y|p − |x|p − p y |x|p−1 sign x − cp |y|p = |x|p 1 +  − 1 − p − cp   x x x and this is nonnegative by (4.27).



By the same method, it can be shown that also the spaces p of sequences and the Sobolev spaces W m,p (D) are p−convex for p ≥ 2 and any m ∈ N0 .

4.3

The Nonlinear Landweber Iteration

As a general criticism towards Tikhonov’s method in the nonlinear case, we mention the disadvantage that the convergence results hold only for the global minima of the Tikhonov functional Jα,δ which is, in general, a non-convex function. In the best of all cases, the global minima can only be computed by iterative methods6 . Therefore, it seems to be natural to solve the nonlinear equation K(x) = y δ directly by an iterative scheme such as Newton-type methods or Landweber methods. In this section, we present the simplest of such an algorithm and follow the presentation of the paper [126] and also the monograph [149]. In this section, let again X and Y be Hilbert spaces and K : X ⊃ D(K) → Y a continuously Fr´echet differentiable mapping from the open domain of definition D(K) and let x ˆ ∈ D(K). First we recall the nonlinear Tikhonov functional Jα,δ (x) = K(x) − y δ 2Y + αx − x ˆ2X ,

x ∈ D(K) ,

for α, δ ≥ 0. This functional is differentiable. Lemma 4.30 Jα,δ is differentiable in every x∗ ∈ D(K) and    Jα,δ (x∗ )h = 2 Re K(x∗ ) − y δ , K  (x∗ )h Y + 2α Re(x∗ − x ˆ, h)X ,

h∈X.

Proof: We have K(x∗ +h) = K(x∗ )+K  (x∗ )h+r(h) with r(h)Y /hX → 0 for h → 0. Using the binomial formula a + b2 = a2 + 2 Re(a, b) + b2 we have Jα,δ (x∗ + h) = K(x∗ + h) − y δ 2Y + αx∗ + h − x ˆ2X

= (K(x∗ ) − y δ ) + (K  (x∗ )h + r(h))2Y + α(x∗ − x ˆ) + h2X   = K(x∗ ) − y δ 2Y + 2 Re K(x∗ ) − y δ , K  (x∗ )h + r(h) Y

+ K  (x∗ )h + r(h)2Y + αx∗ − x ˆ2X + 2α Re(x∗ − x ˆ, h)X

+ αh2X   = Jα,δ (x∗)+2 Re K(x∗)−y δ , K  (x∗ )h Y +2α Re(x∗ − x ˆ, h)X + r˜(h) 6 Even this is only possible in very special cases. Usually, one must be satisfied with critical points of Jα,δ

150

Nonlinear Inverse Problems

with   r˜(h) = 2 Re K(x∗ ) − y δ , r(h) Y + K  (x∗ )h + r(h)2Y + αh2X , thus limh→0 r˜(h)/hX = 0.



 (x∗ )h = 0 for Lemma 4.31 If x∗ ∈ D(K) is a local minimum of Jα,δ then Jα,δ    ∗ ∗ ∗ δ ∗ all h ∈ X, thus K (x ) K(x ) − y + α (x − x ˆ) = 0.

Proof: Let h = 0 fixed. For sufficiently small t > 0 the point x∗ + th belongs to D(K) because D(K) is open. Therefore, for sufficiently small t > 0 the optimality of x∗ implies Jα,δ (x∗ + th) − Jα,δ (x∗ ) /thX ≥ 0, thus 0 ≤

  Jα,δ (x∗ + th) − Jα,δ (x∗ ) − Jα,δ (x∗ )(th) (x∗ )h Jα,δ + . thX hX

 The inequality Jα,δ (x∗ )h ≥ 0 follows as t tends to zero. This implies the assertion since this holds for all h. 

The equation

  ˆ) = 0 K  (x∗ )∗ K(x∗ ) − y δ + α (x∗ − x

(4.28)

reduces to the well known normal equation (2.16) in the case that K is linear and x ˆ = 0 because the derivative is given by K  (x∗ ) = K. However, in general, this equation is nonlinear, and one has to apply an iterative method for solving this equation. Instead of solving the well posed nonlinear equation (4.28), one can directly use an iterative scheme to solve—and regularize—the unregularized  equation K  (x∗ )∗ K(x∗ ) − y δ = 0, which can be written as   x∗ = x∗ − a K  (x∗ )∗ K(x∗ ) − y δ with an arbitrary number a > 0. This is a fixed point equation, and it is natural to solve it iteratively by the fixed point iteration: Fix x ˆ ∈ D(K). Set xδ0 = x ˆ and   xδk+1 = xδk − a K  (xδk )∗ K(xδk ) − y δ , k = 0, 1, 2, . . . This is called the nonlinear Landweber iteration because it reduces to the well known Landweber iteration in the linear case, see Sections 2.3 and 2.6. At the moment, it is not clear that it is well-defined; that is, that all of the iterates lie in the domain of definition D(K). We choose ρ > 0 and a > 0 such that aK  (x)2L(X,Y ) ≤ 1 for all x ∈ B(ˆ x, ρ). δ Then we scale the equation K(x) = y ; that is, we replace it by the equivalent ˜ = √a K and y˜δ = √a y δ . Then K ˜  (x)2 ˜ ≤1 equation K(x) = y˜δ with K for all x ∈ B(ˆ x, ρ), and the Landweber iteration takes the form   ˜ δk ) − y˜δ , k = 0, 1, 2, . . . ˜  (xδk )∗ K(x xδk+1 = xδk − K

L(X,Y )

4.3

The Nonlinear Landweber Iteration

151

Therefore, we assume from now on that a = 1 and K  (x)L(X,Y ) ≤ 1 for all x ∈ B(ˆ x, ρ). The Landweber iteration thus takes the form   xδ0 = x (4.29) ˆ , xδk+1 = xδk − K  (xδk )∗ K(xδk ) − y δ , k = 0, 1, 2, . . . The following, rather strong, assumption is called the “tangential cone condition” and will ensure that the Landweber iteration is well-defined and will also provide convergence. Assumption 4.32 (Tangential Cone Condition) Let B(ˆ x, ρ)⊂D(K) be some ball, K differentiable in B(ˆ x, ρ) with K  (x)L(X,Y ) ≤1  x, ρ), and for all x ∈ B(ˆ x, ρ). Furthermore, let K be Lipschitz continuous on B(ˆ there exists η < 12 with K(˜ x) − K(x) − K  (x)(˜ x − x) Y ≤ η K(˜ x) − K(x)Y (4.30) for all x, x ˜ ∈ B(ˆ x, ρ). We compare this estimate with the estimate K(˜ x) − K(x) − K  (x)(˜ x − x) Y ≤ c ˜ x − x2X , which holds for Lipschitz-continuously differentiable functions (see Lemma A.63). If there exists c with ˜ x − xX ≤ c K(˜ x) − K(x)Y then (4.30) follows. However, such an estimate does not hold for ill-posed problems (see Definition 4.1 and Problem 4.1). The condition (4.30) is very strong but in some cases it can be verified (see Problem 4.7). First we draw some conclusions from this condition. The first one is a basic inequality which will be used quite often in the following. Corollary 4.33 Under Assumption 4.32, the following holds: 1 1 K  (x)(˜ K  (x)(˜ x − x)Y ≤ K(˜ x) − K(x)Y ≤ x − x)Y (4.31) 1+η 1−η for all x, x ˜ ∈ B(ˆ x, ρ). Proof:

The right estimate follows from

K(˜ x) − K(x)Y

≤ ≤

K(˜ x) − K(x) − K  (x)(˜ x − x)Y + K  (x)(˜ x − x)Y  η K(˜ x) − K(x)Y + K (x)(˜ x − x)Y ,

and the left from K(˜ x) − K(x)Y

  = K  (x)(˜ x − x) − K  (x)(˜ x − x) − K(˜ x) + K(x) Y ≥



K  (x)(˜ x − x)Y − η K(˜ x) − K(x)Y .

152

Nonlinear Inverse Problems

With this lemma, we try to justify the notion ”tangential cone condition” and define the convex cones C± (x) in R × X with vertex at (0, x) as   1 K  (x)(x − x ˜) ∈ R × X : r ≥ ˜)Y . C± (x) = (r, x 1±η   ˜ ∈ C+ (x) \ int C− (x) Then (4.31) can be formulated as K(˜ x) − K(x)Y , x for every x, x ˜ ∈ B(ˆ x, ρ). Therefore, for every x ∈ B(ˆ x, ρ) the graph of x ˜ → K(˜ x) − K(x)Y lies between the cones C+ (x) and C− (x) which are build by the tangent space at x. We draw two conclusions from this corollary. Under Assumption 4.32, the reverse of Theorem 4.3 holds, and there exists a unique minimum norm solution with respect to x ˆ. Definition 4.34 Let K : X ⊃ D(K) → Y be a nonlinear mapping, y ∗ ∈ Y , and x ˆ ∈ X. A point x∗ ∈ D(K) with K(x∗ ) = y ∗ is called minimum norm ˆX ≤ x − x ˆX for all solutions x ∈ D(K) solution with respect to x ˆ if x∗ − x ∗ of K(x) = y . Lemma 4.35 Let Assumption 4.32 hold for some ball B(ˆ x, ρ) and let x∗ ∈ B(ˆ x, ρ) with K(x∗ ) = y ∗ . (a) Let the linear equation K  (x∗ )h = 0 be ill-posed in the sense of Definition 4.1. Then the nonlinear equation K(x) = y is locally ill-posed in x∗ . (b) x∗is a minimum norm solution with respect to x ˆ if, and only if, x∗ − x ˆ⊥   ∗ N K (x ) (c) There exists a unique minimum norm solution x† of K(x) = y ∗ with respect to x ˆ.   ˆX . Since the linear equation is ill-posed Proof: (a) Let r ∈ 0, ρ − x∗ − x there exists a sequence hn ∈ X with hn X = 1 and K  (x∗ )hn Y → 0 for ˜ = x∗ and x = xn in (4.31) it follows that n → ∞. We set xn = x∗ + rhn . For x ∗ ∗ K(xn ) → K(x ) and xn − x X = r. (b) For the characterization of a minimum norm solution it suffices to compare x∗ − x ˆX with ˜ x−x ˆX for solutions x ˜ ∈ B(ˆ x, ρ) of K(˜ x) = y ∗ . The estimate ˜ ∈ B(ˆ x, ρ) is a solution of K(˜ x) = y ∗ if, and only (4.31) for x = x∗ implies that x  x) = 0; that is, if x∗ −˜ x ∈ N K  (x∗ ) . The point x∗ is a minimum if, K  (x∗ )(x∗ −˜ ˆX ≤ ˜ x−x ˆX for all norm solution with respect to x ˆ if, and only if, x∗ − x  ˜ with x∗ − x ˜ ∈ N K  (x∗ ) . solutions x ˜ ∈ B(ˆ x, ρ) of K(˜ x) = y ∗ ; that is, for all x ˜ by z shows that x∗ is a minimum norm solution Replacing x∗ − x  with respect ∗ ∗ − x ˆ in N K (x ) which to x ˆ if, and only if, 0 is the best approximation of x  ∗  ∗ ˆ ⊥ N K (x ) . is equivalent to x − x   ∗  (c) The subspace N K (x ) is closed. Therefore, there exists a unique best  ˆ in N K  (x∗ ) .7 Then x† = x∗ − p ∈ x∗ + N K  (x∗ ) approximation p of x∗ − x 7 This

follows from a general theorem on best approximations in Hilbert spaces.

4.3

The Nonlinear Landweber Iteration

153

  is the best approximation at x ˆ in x∗ + N K  (x∗ ) . Also, K(x† ) = y ∗ because   ˆX ≤ x∗ − x ˆX < ρ, thus x† ∈ x∗ − x† = p ∈ N K  (x∗ ) . Finally, x† − x B(ˆ x, ρ).  The following result proves that the Landweber iteration is well-defined, and it motivates a stopping rule. Theorem 4.36 Let Assumption 4.32 hold. Let y δ − y ∗ Y ≤ δ and x∗ ∈ B(ˆ x, ρ/2) with K(x∗ ) = y ∗ . Furthermore, let r > 2(1+η) 1−2η and assume that δ ∗ there exists k∗ ∈ N with xk ∈ B(x , ρ/2) for all k = 0, . . . , k∗ − 1 (where xδk are defined by (4.29)) such that K(xδk ) − y δ Y ≥ rδ

for all k = 0, . . . , k∗ − 1 .

(4.32)

Then the following holds. (a) xδk+1 − x∗ X ≤ xδk − x∗ X for all k = 0, . . . , k∗ − 1. In particular, xδk − x∗ X ≤ xδ0 − x∗ X = ˆ x − x∗ X < ρ/2 for all k = 0, . . . , k∗ . δ ∗ x, ρ) for k = 0, . . . , k∗ . Therefore, all xk belong to B(x , ρ/2) ⊂ B(ˆ (b) For all ∈ {0, . . . , k∗ − 1} it holds that   k∗ −1 2(1 + η)

1 − 2η − K(xδk ) − y δ 2Y ≤ xδ − x∗ 2X − xδk∗ − x∗ 2X . r k= (4.33) (c) If δ = 0 then (4.32) holds for all k ∈ N, and (we write xk instead of x0k ) ∞

K(xk ) − y ∗ 2Y ≤

k=0

ˆ x − x∗ 2X . 1 − 2η

(4.34)

In particular, K(xk ) converges to y ∗ . Proof:

Since K  (xδk )∗ L(Y,X) ≤ 1 we have for k = 0, . . . , k∗ − 1

xδk+1 − x∗ 2X − xδk − x∗ 2X = = ≤ = ≤ ≤ = = ≤

2 Re(xδk − x∗ , xδk+1 − xδk )X + xδk+1 − xδk 2X 2   2 Re xδk − x∗ , K  (xδk )∗ (y δ − K(xδk )) X + K  (xδk )∗ (y δ − K(xδk )) X   2 Re K  (xδk )(xδk − x∗ ), y δ − K(xδk ) X + y δ − K(xδk )2Y   2 Re K  (xδk )(xδk − x∗ ) − K(xδk ) + y δ , y δ − K(xδk ) X − y δ − K(xδk )2Y   2η K(xδk ) − y ∗ Y + 2δ − K(xδk ) − y δ Y K(xδk ) − y δ Y   (2η − 1) K(xδk ) − y δ Y + 2δ + 2ηδ K(xδk ) − y δ Y   2δ(1 + η) − (1 − 2η) K(xδk ) − y δ Y K(xδk ) − y δ Y (4.35)  1 2(1 + η)δr − (1 − 2η)r K(xδk ) − y δ Y K(xδk ) − y δ Y r 2(1 + η) − (1 − 2η)r K(xδk ) − y δ 2Y r

154

Nonlinear Inverse Problems

where we used (4.32) in the last estimate. This implies (a), because 2(1 + η) − (1 − 2η)r < 0 by the choice of r.   With α = (1 − 2η)r − 2(1 + η) /r > 0 we just have shown the estimate α K(xδk ) − y δ 2Y ≤ xδk − x∗ 2X − xδk+1 − x∗ 2X ,

k = 0, . . . , k∗ − 1 ,

and thus for ∈ {0, . . . , k∗ − 1}: α

k

∗ −1

K(xδk ) − y δ 2Y ≤ xδ − x∗ 2X − xδk∗ − x∗ 2X .

k=

(c) We consider the above estimate up to line (4.35) for δ = 0; that is, (1 − 2η) K(xk ) − y ∗ 2Y ≤ xk − x∗ 2X − xk+1 − x∗ 2X , and thus for all m ∈ N (1 − 2η)

m−1

K(xk ) − y ∗ 2Y ≤ x0 − x∗ 2X − xm − x∗ 2X ≤ x0 − x∗ 2X .

k=0

Therefore the series converges which yields the desired estimate.



Let now δ > 0 and again r > 2(1+η) condition (4.32) can not hold for 1−2η . The   every k∗ . Indeed, otherwise the sequence xk −x∗ X would be a monotonically decreasing and bounded sequence, thus convergent. From (4.33) for = k∗ − 1 it would follows that   2(1 + η) 1 − 2η − K(xδk∗ −1 )−y δ 2Y ≤ xδk∗ −1 −x∗ 2X − xδk∗ −x∗ 2X −→ 0 r for k∗ → ∞, a contradiction. Therefore, the following stopping rule is welldefined. Stopping rule: Let δ > 0 and r > 2(1+η) 1−2η . We define k∗ = k∗ (δ) ∈ N as the uniquely determined number such that K(xδk ) − y δ Y ≥ rδ > K(xδk∗ ) − y δ Y

for all k = 0, . . . , k∗ − 1 . (4.36)

With this stopping rule, we can show convergence as δ → 0. First we note that for any fixed δ > 0 and k ≤ k∗ (δ), the mapping y δ → xδk is continuous. Indeed, in every of the first k steps of the algorithm only continuous operations of the Landweber iteration are performed. We first consider the case of no noise; that is, δ = 0. Theorem 4.37 Let Assumption 4.32 hold and let ρ > 0 and x∗ ∈ B(ˆ x, ρ/2) with K(x∗ ) = y ∗ . The Landweber iteration is well defined for δ = 0 (that is, the iterates xk belong to D(K)) and the sequence (xk ) converges to a solution x, ρ) of K(x) = y ∗ . x ˜ ∈ B[x∗ , ρ/2] ⊂ B(ˆ

4.3

The Nonlinear Landweber Iteration

155

Proof: We have seen already in Theorem 4.36 that xk ∈ B(x∗ , ρ/2) for all k. We now show that (xk ) is a Cauchy sequence. Let ≤ m be fixed. Determine k with ≤ k ≤ m and

K(xk ) − y ∗ Y ≤ K(xi ) − y ∗ Y

for all ≤ i ≤ m .

Because of x − xm 2 ≤ 2x − xk 2 + 2xk − xm 2 we estimate both terms separately and use the formula u − v2 + u2 − v2 = 2 Re(u, u − v) and the left estimate of (4.31).

xk − xm 2X + xk − x∗ 2X − xm − x∗ 2X m−1

    = 2 Re xk − x∗ , xk − xm X = 2 Re xk − x∗ , xi − xi+1 X i=k

=2

m−1



Re xk − x∗ , K  (xi )∗ (K(xi ) − y ∗ )

i=k

=2

m−1

i=k

≤2

m−1

 X

  Re K  (xi )(xk − x∗ ), K(xi ) − y ∗ Y K(xi ) − y ∗ Y K  (xi )(xk − x∗ )Y

i=k

≤2

m−1

  K(xi ) − y ∗ Y K  (xi )(xi − x∗ )Y + K  (xi )(xi − xk )Y

i=k

≤ 2(1 + η)

m−1

  K(xi ) − y ∗ Y K(xi ) − K(x∗ )Y + K(xi ) − K(xk )Y

i=k

≤ 2(1+η)

m−1

  K(xi )−y ∗ Y K(xi )−y ∗)Y + K(xi )−y ∗ Y + y ∗ − K(xk )Y

i=k

≤ 2(1+η)

m−1

  K(xi )−y ∗ Y K(xi )−y ∗ Y +K(xi )−y ∗ Y + y ∗ − K(xi )Y

i=k

= 6(1 + η)

m−1

K(xi ) − y ∗ 2Y .

i=k

Analogously, we estimate

xk − x 2X + xk − x∗ 2X − x − x∗ 2X ≤ 6(1 + η)

k−1

i=

K(xi ) − y ∗ 2Y

156

Nonlinear Inverse Problems

and thus x − xm 2X

≤

2x − xk 2X + 2xk − xm 2X

≤

2x − x∗ 2X − 2xk − x∗ 2X + 12(1 + η)

k−1

K(xi ) − y ∗ 2Y

i=

+2xm − x∗ 2X − 2xk − x∗ 2X + 12(1 + η)

m−1

K(xi ) − y ∗ 2Y

i=k

=

2x − x∗ 2X + 2xm − x∗ 2X − 4xk − x∗ 2X + 12(1 + η)

m−1

K(xi ) − y ∗ 2Y .

i=

 

∞ Since both, the sequence x − x∗ X and also the series i=0 K(xi ) − y ∗ 2Y , converge by Theorem 4.36 we conclude that also x − xm 2X converges to zero as → ∞. Therefore, (xm ) is a Cauchy sequence and thus convergent; that is, xm → x ˜ for some x ˜ ∈ B[x∗ , ρ/2]. The continuity of K implies that ∗ K(˜ x) = y .  Now we prove convergence of xδk∗ (δ) as δ → 0 if k∗ (δ) is determined by the stopping rule. Theorem 4.38 Let Assumption 4.32 hold and let ρ > 0 and x∗ ∈ B(ˆ x, ρ/2) . Then the sequence stops by the stopping rule with K(x∗ ) = y ∗ and r > 2(1+η) 1−2η δ ∗ (4.36) and defines k∗ (δ) for all δ > 0. Then lim xk∗ (δ) = x ˜ where x ˜ ∈ B[x , ρ/2] δ→0

is the limit of the sequence for δ = 0 – which exists by the previous theorem and is a solution of K(˜ x) = y ∗ . Proof: Let (xk ) be the sequence of the Landweber iteration for δ = 0; that ˜ = limk→∞ xk . Furthermore, let (δn ) be a sequence which is, xk = x0k , and x converges to zero. Set for abbreviation kn = k∗ (δn ). We distinguish between two cases. Case 1: The sequence (kn ) of natural numbers has finite accumulation points. Let k ∈ N be the smallest accumulation point. Then there exists I ⊂ N of infinite cardinality such that kn = k for all n ∈ I. By the definition of kn , we have K(xδkn ) − y δn Y < r δn

for all n ∈ I .

Since for this fixed k, the iteration xδkn depends continuously on y δn , and since y δn converges to y ∗ , we conclude for n ∈ I, n → ∞, that xδkn → xk (this is the k-th iteration of the sequence for δ = 0) and K(xk ) = y ∗ . Landweber’s iteration for δ = 0 implies that then xm = xk for all m ≥ k; that is, the sequence is ˜ for all m ≥ k. The same argument constant for m ≥ k. In particular, xm = x x. holds for any other accumulation point k˜ ≥ k and thus lim xδkn∗ (δn ) = xk˜ =˜ n→∞,n∈I˜

4.3

The Nonlinear Landweber Iteration

157

Case 2: The sequence (kn )n tends to infinity. Without loss of generality, we assume that kn converges monotonically to infinity. Let n > m. We apply Theorem 4.36 for x ˜ instead of x∗ and get xδknn − x ˜X

≤

xδknn−1 − x ˜X ≤ · · · ≤ xδknm − x ˜X

≤

xδknm − xkm X + xkm − x ˜X .

Let ε > 0. Choose m such that xkm − x ˜X ≤ ε/2. For this fixed m, the sequence xδknm converges to xkm for n → ∞. Therefore, we can find n0 with ˜X ≤ ε. xδknm − xkm X ≤ ε/2 for all n ≥ n0 . For these n ≥ n0 , we have xδknn − x Therefore, we have shown convergence of every sequence δn → 0 to the same limit x ˜. This ends the proof.  Before we prove a result on the order of convergence, we formulate a condition under which the Landweber iteration converges to the minimum norm solution. x, ρ/2) be the minimum Lemma 4.39 Let Assumption 4.32 hold and let x∗ ∈ B(ˆ ˆ. (It exists and is unique by norm solution of K(x) = y ∗ with respect to x  x, ρ/2). Lemma 4.35.) Let, furthermore, N K  (x∗ ) ⊂ N K  (x) for all x ∈ B(ˆ  n  of the Landweber iteration converge to Then the sequences (xk ) and xδk(δ n) n the minimum norm solution x∗ of K(x) = y ∗ . Proof: Convergence of the sequences has been shown in Theorems 4.37 and 4.38 already. Let x ˜ = limk→∞ xk . Then also x ˜ = limδ→0 xδk(δ) by Theorem 4.38.   We show by induction that xk − x∗ ⊥ N K  (x∗ ) for all k = 0, 1, . . .. For k = 0 ˆ, the minimum property of x∗ , and part this is true because x0 = x  (b) of  ∗ (x ) . Then Lemma 4.35. Let the assertion be true for k and let z ∈ N K    z ∈ N K (xk ) , thus   (xk+1 − x∗ , z)X = (xk − x∗ , z)X + K(xk ) − y ∗ , K  (xk )z Y = 0 .   ∗  Since the orthogonal complement of N K (x ) is closedwe conclude that also   ˜ − x∗ ∈ N K  (x∗ ) because of (4.31) for x ˜ − x∗ ⊥ N K  (x∗ ) . Furthermore, x ˜ = x∗ .  x = x∗ . Therefore, x We will now prove a rate of convergence. This part will be a bit more technical. It is not surprising that we need the source condition of Assumption 4.11 just as in the linear case or for the nonlinear Tikhonov regularization. Unfortunately, Assumptions 4.32 and 4.11 are not sufficient, and we have to strengthen them. x, ρ) be the minimum Assumption 4.40 Let Assumption 4.32 hold and x∗ ∈ B(ˆ ˆ. (It exists by Lemma 4.35). norm solution of K(x) = y ∗ with respect to x Furthermore, let K  (x∗ ) compact and there exists C > 0 and a family {Rx : x ∈ B(ˆ x, ρ)} of linear bounded operators Rx : Y → Y with K  (x) = Rx K  (x∗ )

and

Rx − IL(Y,Y ) ≤ C x − x∗ X

for x ∈ B(ˆ x, ρ) .

158

Nonlinear Inverse Problems

In the linear case, both conditions are satisfied for Rx = I because K  (x) = K for all x. Under this additional assumption, we can partially sharpen the tangential cone condition of Assumption 4.32.     Lemma 4.41 Under Assumption 4.40, we have N K  (x∗ ) ⊂ N K  (x) for all x ∈ B(ˆ x, ρ) and K(x) − K(x∗ ) − K  (x∗ )(x − x∗ ) Y K(x) − K(x∗ ) − K  (x∗ )(x − x∗ ) Y

C x − x∗ X K  (x∗ )(x − x∗ )Y , 2 C ≤ x − x∗ X K(x) − K(x∗ )Y 2 − ρC ≤

for all x ∈ B(ˆ x, ρ). Therefore, if we choose ρ such that η := condition (4.30) is satisfied for x = x∗ . Proof:

Cρ 2−ρC

0

1 =



   K  tx + (1 − t)x∗ − K  (x∗ ) (x − x∗ ) dt ,

0

thus K(x) − K(x∗ ) − K  (x∗ )(x − x∗ )Y 1      K tx + (1 − t)x∗ − K  (x∗ ) (x − x∗ ) dt ≤ Y 0

1 ≤

Rtx+(1−t)x∗ − IL(Y,Y ) dt K  (x∗ )(x − x∗ )Y

0

1

t dt x − x∗ X K  (x∗ )(x − x∗ )Y

≤

C

=

C x − x∗ X K  (x∗ )(x − x∗ )Y 2

0

which proves (4.37). We continue and estimate K(x) − K(x∗ ) − K  (x∗ )(x − x∗ ) Y C x − x∗ X K  (x∗ )(x − x∗ ) + K(x∗ ) − K(x) Y ≤ 2 C x − x∗ X K(x∗ ) − K(x)Y + 2

(4.38)

< 1/2 then

The first assertion is obvious. For the second, we write we K(x) − K(x∗ ) − K  (x∗ )(x − x∗ )  1   d  ∗  ∗ ∗ K tx + (1 − t)x − K (x )(x − x ) dt = dt

(4.37)

4.3

The Nonlinear Landweber Iteration

159

and thus   2−Cx−x∗ X K(x)−K(x∗ )−K  (x∗)(x−x∗) Y ≤C x−x∗ X K(x∗ )−K(x)Y . This proves (4.38) because x − x∗ X ≤ ρ.



x, ρ/2)), we know from Lemma 4.39 Under this Assumption 4.40 (and x∗ ∈ B(ˆ that xδk(δ) converges to x∗ as δ → 0. For the proof of the order of convergence, we need the following elementary estimates: Lemma 4.42 1 − (1 − x2 )k x

≤

x (1 − x2 )k

≤

x2 (1 − x2 )k

≤

√

for all k ∈ N and x ∈ [0, 1] ,

k

1 for all k ∈ N and x ∈ [0, 1] , k+1 1 for all k ∈ N and x ∈ [0, 1] , k+1 √

(4.39) (4.40) (4.41)

for every α ∈ (0, 1] there exists c(α) > 0 with k−1

(j + 1)−α (k − j)−3/2 ≤

j=0

c(α) (k + 1)α

for all k ∈ N .

(4.42)

√ Proof: To show (4.39) let first 0 < x ≤ 1/ k. Bernoulli’s inequality implies √ √ 2 k 2 ) ) (1−x2 )k ≥ 1−kx2 , thus 1−(1−x ≤ 1−(1−kx = kx ≤ k. Let now x ≥ 1/ k. x x √ 2 k ) Then 1−(1−x ≤ x1 ≤ k. For x = 0 the estimate follows by l’Hospital’s rule. x The proofs of estimates (4.40) and (4.41) are elementary and left to the reader (see also the proof of Theorem 2.8). Proof of (4.42): Define the function f (x) = (x + 1)−α (k − x)−3/2 for −1 < x < k. By the explicit computation of the second derivative, one obtains that f  (x) ≥ 0 for all −1 < x < k.8 Taylor’s formula yields for j < k 1  f (zj,x )(x − j)2 ≥ f (j) + f  (j)(x − j) , 2

f (x) = f (j) + f  (j)(x − j) +

for j − 1/2 ≤ x ≤ j + 1/2 with some intermediate point zj,x , thus j+1/2 

j+1/2 



f (x) dx ≥ j−1/2

8 Therefore,

f is convex.

j−1/2

 f (j) + f  (j)(x − j) dx = f (j) ,

160

Nonlinear Inverse Problems

because the integral

j+1/2  j−1/2

(x − j) dx vanishes by the symmetry of the intervals.

Therefore, k−1

−1/2

(j + 1)

−3/2

(k − j)

=

j=0

k−1

f (j) ≤

j=0

k−1

f (x) dx =

j=0 j−1/2

k/2 =

k−1/2 

j+1/2 

f (x) dx −1/2

k−1/2 

f (x) dx + −1/2

f (x) dx . k/2

For the first integral, we have k/2 −1/2

 3/2 k/2 2 f (x) dx ≤ (x + 1)−α dx . k −1/2

If α < 1 then k/2 −1/2

1−α  3/2  2 k 1 c(α) +1 f (x) dx ≤ ≤ c˜(α) k −α−1/2 ≤ . k 1−α 2 (k + 1)α

If α = 1 then k/2 f (x) dx ≤ −1/2

 3/2 2 c(1) . ln(k/2) ≤ k k+1

For the remaining integral, we estimate ≤

 α k−1/2  α  k−1/2 2 2 (k − x)−3/2 dx = 2 (k − x)−1/2 k/2 k k

≤

 α √ 2 2 2. k

k−1/2 

f (x) dx k/2

k/2

This finishes the proof.



(4.39), (4.40), and (4.41) directly imply Corollary 4.43 Let A : X → Y be a linear and compact operator with adjoint A∗ and let AL(X,Y ) ≤ 1. Then (a) (I − A∗ A)k A∗ L(Y,X) ≤ (k + 1)−1/2 ,

4.3

The Nonlinear Landweber Iteration

161

(b) (I − A A∗ )k A A∗ L(Y,Y ) ≤ (k + 1)−1 , k−1 √ ∗ j ∗ (c) (I − A A) A ≤ k. j=0 L(Y,X)

Proof: Let {μi , xi , yi : i ∈ I} be a singular system for A, see Appendix A.6, Theorem A.57. For y = y0 + i∈I αi yi ∈ Y with A∗ y0 = 0, we have

 2 (I − A∗ A)k A∗ y2X = (1 − μ2i )k μi αi2 ≤ i∈I

1 2 1 y2Y αi ≤ k+1 k+1 i∈I

because of (4.40). The part (b) follows in the way from (4.41). Finally, we observe that k−1

(I − A∗ A)j A∗ y =

j=0

k−1

j=0 i∈I

1 − (1 − μ2 )k   i (1 − μ2i )j μi αi yi = αi yi μi i∈I

and thus by (4.39) 2 k−1



(I − A∗ A)j A∗ y ≤ k αi2 ≤ k y2Y . j=0 i∈I Y

 Now we are able to prove the main theorem on the order of convergence. The proof is lengthy. Theorem 4.44 Let Assumption 4.40 hold, and let x∗ ∈ B(ˆ x, ρ) be the minimum ˆ. Define the constant norm solution of K(x) = y ∗ with respect to x   1 − 2η c∗ = 2 1 + η (5 − 4η) where η < 1/2 is the constant of Assumption 4.32. Furthermore, let the source ˆ = K  (x∗ )∗ w and condition hold; that is, there exists w ∈ Y such that x∗ − x  9 C c2∗ max c(1/2), c(1) wY ≤ 1, 2 where C is the constant of Assumption 4.40 and c(α) are the constants of (4.42) for α = 1/2 and α = 1. Then there exists a constant c > 0 such that 1/2

xδk∗ (δ) − x∗ X ≤ c wY δ 1/2 ,

K(xδk∗ (δ) ) − y ∗ Y ≤ (1 + r) δ . (4.43)

Again, k∗ (δ) is the iteration index determined by the stopping rule (4.36).

162

Nonlinear Inverse Problems

Proof: The second estimate follows immediately from the property of k∗ (δ) and the triangle inequality. Set for abbreviation ek = xδk − x∗ and A = K  (x∗ ). Then K  (xδk ) = Rxδk A because of Assumption 4.40 and thus ek+1

= = =

=

  ek − K  (xδk )∗ K(xδk ) − y δ     ek − A∗ K(xδk ) − y δ + A∗ (I − Rx∗ δ ) K(xδk ) − y δ k   (I − A∗ A)ek − A∗ K(xδk ) − y δ − K  (x∗ )ek   + A∗ (I − Rx∗ δ ) K(xδk ) − y δ k   (I − A∗ A)ek + A∗ (y δ − y ∗ ) − A∗ K(xδk ) − K(x∗ ) − K  (x∗ )(xδk − x∗ )   + A∗ (I − Rx∗ δ ) K(xδk ) − y δ k

=

(I − A∗ A)ek + A∗ (y δ − y ∗ ) + A∗ zk ,

k = 0, . . . , k∗ (δ) − 1 ,

(4.44)

with     zk = (I − Rx∗ δ ) K(xδk ) − y δ − K(xδk ) − K(x∗ ) − K  (x∗ )(xδk − x∗ ) . (4.45) k

This is a recursion formula for ek . It is solved in terms of y δ − y ∗ and zj by ek

=

(I − A∗ A)k e0 +

k−1

  (I − A∗ A)j A∗ (y δ − y ∗ ) + zk−j−1

j=0

⎤

⎡

k−1

= −(I − A∗ A)k A∗ w + ⎣

(I − A∗ A)j A∗ ⎦(y δ − y ∗ )

j=0

+

k−1

(I − A∗ A)j A∗ zk−j−1

(4.46)

j=0

for k = 0, . . . , k∗ (δ). (Proof by induction with respect to k.) Furthermore, because A(I − A∗ A)j = (I − A A∗ )j A,

4.3

Aek

The Nonlinear Landweber Iteration

= −(I − A A∗ )k A A∗ w +

163

k−1

(I − A A∗ )j A A∗ (y δ − y ∗ )

j=0

+

k−1

(I − A A∗ )j A A∗ zk−j−1

j=0

= −(I − A A∗ )k A A∗ w −

k−1

  (I − A A∗ )j+1 − [(I − A A∗ )j (y δ − y ∗ )

j=0

+

k−1

(I − A A∗ )j A A∗ zk−j−1

j=0

  = −(I − A A∗ )k A A∗ w + I − (I − A A∗ )k (y δ − y ∗ ) +

k−1

(I − A A∗ )j A A∗ zk−j−1

(4.47)

j=0

for k = 0, . . . , k∗ (δ). Now we consider zk from (4.45) for k ∈ {0, . . . , k∗ (δ) − 1} and estimate both terms of zk separately. Since r > 2 and K(xδk ) − y δ Y ≥ rδ > 2δ we have K(xδk )−y δ Y ≤ K(xδk )−y ∗ Y + δ ≤ K(xδk ) − y ∗ Y +

1 K(xδk ) − y δ Y 2

and thus K(xδk ) − y δ Y ≤ 2 K(xδk ) − y ∗ Y . With (4.31) and because η < 1/2 we have   (I −R∗ δ ) K(xδk )−y δ ≤I −R∗ δ L(Y,Y ) K(xδk )−y δ Y ≤ 4C ek X Aek Y . x x Y k

k

For second term of zk we use (4.37) for x = xδk : K(xδk ) − K(x∗ ) − K  (x∗ )(xδk − x∗ ) ≤ C ek X Aek Y Y 2 and thus zk Y ≤ 92 C ek X Aek Y . Now we estimate δ from above. Because of the stopping rule (4.36) and (4.31) we have for k = 0, . . . , k∗ (δ) − 1 r δ ≤ K(xδk ) − y δ Y ≤ K(xδk ) − K(x∗ )Y + δ ≤ and thus, because r − 1 > δ ≤

2(1+η) 1−2η

−1=

1 Aek Y + δ , 1−η

1+4η 1−2η ,

1 − 2η Aek Y , (1 + 4η)(1 − η)

k = 0, . . . , k∗ (δ) − 1 .

(4.48)

164

Nonlinear Inverse Problems

Now we show ej X ≤ c∗ wY √

1 , j+1

Aej Y ≤ c∗ wY

1 , j+1

(4.49)

and

2(1 − 2η) wY . (4.50) η(5 − 4η) k + 1 for j = 0, . . . , k and every k ≤ k∗ (δ) − 1. We show these three estimates by induction with respect to k < k∗ (δ). These are true for k = 0 because 2(1−2η) 1−2η ≤ η(5−4η) . e0 X ≤ wY and Ae0 Y ≤ wY and by (4.48) and (1+4η)(1−η) Let now (4.49) and (4.50) hold for k − 1. From the above representation (4.46) of ek , parts (a) and (c) of Corollary 4.43, the assumption of induction, and (4.42) for α = 1/2 we conclude that δ ≤

ek X

≤

k−1 √ 1 w 9

√ Y + kδ + C √ ek−j−1 X Aek−j−1 Y 2 j=0 j + 1 k+1

≤

k−1

√ 1 1 w 9 1 √ Y + k δ + C c2∗ w2Y √ √ 2 k − j j + 1 k − j k+1 j=0

≤ ≤

√ w c(1/2) 9 √ Y + k δ + C c2∗ w2Y √ 2 k+1 k+1 √ wY k + 1δ + 2 √ k+1

by assumption on wY . Analogously, it follows for Aek Y with (4.47) and part (b) of Corollary 4.43 and (4.42) for α = 1 with constant c(1): Aek Y

≤

k−1 9 1 wY + δ + C ek−j−1 X Aek−j−1 Y k+1 2 j=0 j + 1

≤

k−1

1 9 1 wY 1 √ + δ + C c2∗ w2Y k+1 2 j+1 k−j k−j j=0

9 c(1) wY + δ + C c2∗ w2Y k+1 2 k+1 wY ≤ δ + 2 . k+1 Next we substitute this estimate of Aek Y into the estimate (4.48) for δ and obtain   1 − 2η 1 − 2η wY Aek Y ≤ δ ≤ δ+2 . (1 + 4η)(1 − η) (1 + 4η)(1 − η) k+1 ≤

Since 1 −

1−2η (1+4η)(1−η)

=

η(5−4η) (1+4η)(1−η)

δ ≤

it follows estimate (4.50) for k; that is,

2(1 − 2η) wY . η(5 − 4η) k + 1

4.3

The Nonlinear Landweber Iteration

165

We substitute this into the estimates of ek X and Aek Y : ek X

≤

Aek Y

≤

  w w 1 − 2η √ Y = c∗ √ Y , 2 1+ η(5 − 4η) k+1 k+1   wY wY 1 − 2η = c∗ . 2 1+ η(5 − 4η) k + 1 k+1

This shows (4.49) for k. Therefore, (4.49) and (4.50) are proven for all 0 ≤ j, k < k∗ (δ). In the case k∗ (δ) ≥ 1, we take k = k∗ (δ) − 1 in (4.50) and obtain k∗ (δ) ≤

2(1 − 2η) wY . η(5 − 4η) δ

(4.51)

From the estimate K(xδk ) − y δ Y ≤ 2K(xδk ) − K(x∗ )Y ≤ 4 Aek Y for k < k∗ (δ) it follows that K(xδk ) − y δ Y ≤ 4 c∗ wY

1 , k+1

2 1−η

Aek Y ≤

k = 0, . . . , k∗ (δ) − 1 .

(4.52)

Now we come to the final part of the proof. We write ek from (4.46) in der form ⎤ ⎡ k−1

ek = −(I − A∗ A)k A∗ w + ⎣ (I − A∗ A)j A∗ ⎦(y δ − y ∗ ) j=0

+

k−1

(I − A∗ A)j A∗ zk−j−1

j=0

⎤

⎡

k−1

= −A∗ (I − AA∗ )k w + ⎣

(I − A∗ A)j A∗ ⎦(y δ − y ∗ )

j=0

+ A∗

k−1

j=0

(I − AA∗ )j zk−j−1 ⎤

⎡

k−1

(I − A∗ A)j A∗ ⎦(y δ − y ∗ )

= A∗ wk + ⎣

j=0

for k = 0, . . . , k∗ (δ) with wk = −(I − AA∗ )k w +

k−1

j=0

(I − AA∗ )j zk−j−1 .

(4.53)

166

Nonlinear Inverse Problems

Since (I − AA∗ )j L(Y,Y ) ≤ 1 for all j, we conclude wk Y

≤

k−1

wY +

zk−j−1 Y ≤ wY +

j=0

≤

≤

k−1 9

C ek−j−1 X Aek−j−1 Y 2 j=0

k−1 ∞



1 1 1 9 9 √ C c2∗ w2Y ≤ wY + C c2∗ w2Y 2 k−j k−j 2 j 3/2 j=0 j=1 ⎡ ⎤ ∞

1 9 2 ⎦ C c∗ wY ⎣1 + 2 j 3/2 j=1

wY +

for all k = 0, . . . , k∗ (δ). Here we used without loss of generality that wY ≤ 1. The bound on the right-hand side is independent of δ. Similarly, we estimate AA∗ wk Y . Indeed, with ⎤ ⎡ k−1

(I − A∗ A)j A∗ ⎦(y δ − y ∗ ) AA∗ wk = Aek − ⎣A j=0

from (4.53) and

A

k−1

(I − A∗ A)j A∗

=

j=0

k−1

(I − AA∗ )j AA∗

j=0

=

−

k−1

k−1

j=0

j=0

(I − AA∗ )j+1 +

(I − AA∗ )j = I − (I − AA∗ )k

we estimate AA∗ wk Y

≤

k−1

∗ j ∗ A Aek Y + (I − A A) A j=0

≤

Aek Y

L(Y,Y )

y δ − y ∗ Y

+ I − (I − AA∗ )k Y δ ≤ Aek Y + δ

for k = 0, . . . , k∗ (δ). For k = k∗ (δ) the estimate (4.31) and the stopping rule imply that AA∗ wk∗ (δ) Y

≤ ≤

(1 + η)K(xδk∗ (δ) ) − y ∗ Y + δ     (1 + η) K(xδk∗ (δ) ) − y δ Y + δ + δ ≤ (1 + η)(1 + r) + 1 δ .

This implies A∗ wk∗ (δ) 2X

    = A∗ wk∗ (δ) , A∗ wk∗ (δ) X = AA∗ wk∗ (δ) , wk∗ (δ) Y ≤ AA∗ wk∗ (δ) Y wk∗ (δ) Y ≤ c1 δ wY

4.4

Problems

167

with some constant c1 which is independent of δ and w. Now we go back to (4.53) and get, using part (c) of Corollary 4.43 and (4.51) (in the case k∗ (δ) ≥ 1, otherwise directly) k∗ (δ)−1 +

∗ j ∗ ek∗ (δ) X ≤ c1 δ wY + (I − A A) A δ ≤ ≤

+ +

c1 δ wY +

+ ,

c1 δ wY +

This, finally, ends the proof.

j=0

L(X,Y )

k∗ (δ) δ √ √ 2(1 − 2η) 1/2 1/2 wY δ ≤ c wY δ. η(5 − 4η)



The Landweber iteration is only one member of the large class of iterative regularization methods. In particular, Newton-type methods, combined with various forms of regularization, have been investigated in the past and are subject of current research. We refer to the monograph [149] of Kaltenbacher, Neubauer, and Scherzer.

4.4

Problems

4.1 Let K : X ⊃ D(K) → K be a mapping with ˜ x −xX ≤ cK(˜ x)−K(x)Y for all x, x ˜ ∈ B(x∗ , ρ) where c is independent of x ˜ and x. Set y ∗ = K(x∗ ). Show that the equation K(x) = y ∗ is not locally ill-posed in the sense of Definition 4.1. 4.2 Define the integral operator (auto-convolution) K from L2 (0, 1) into itself by t K(x)(t) = x(t − s) x(s) ds , t ∈ (0, 1) . 0

(a) Show that K is well-defined; that is, that K(x) ∈ L2 (0, 1) for every x ∈ L2 (0, 1). Remark: K is even well-defined as a mapping from L2 (0, 1) into C[0, 1]. One try to prove this. (b) Show that K(x) = y is locally ill-posed in the sense of Definition 4.1 in every x∗ ∈ L2 (0, 1) with x∗ (t) ≥ 0 for almost all t ∈ (0, 1). Hint: For any r > 0 and n ∈ N define xn ∈ L2 (0, 1) by  x∗ (t) √ , 0 < t < 1 − 1/n , xn (t) = ∗ x (t) + r n , 1 − 1/n < t < 1 . 4.3 Show that for every t ∈ [0, 2] there exists ct > 0 such that μt ≤ ct αt/2−1 +α

μ2

for all μ, α > 0 .

168

Nonlinear Inverse Problems

4.4 Show that K(xα(δ),δ ) − y δ + αwY = O(δ) as δ → 0 for the choice c− δ 2/(σ+1) ≤ α(δ) ≤ c+ δ 2/(σ+1) and xα,δ and w as in Theorem 4.15. 4.5 (a) Show that u∞ ≤ u L2 (0,1) for all u ∈ H 1 (0, 1) with u(0) = 0. (b) Show that H 2 (0, 1) is compactly imbedded in C[0, 1]. (c) Show that the solution u ∈ H 2 (0, 1) of u = −g in (0, 1) and u(0) = 1 u(1) = 0 is given by u(t) = 0 G(t, s)g(s)ds where G(t, s) is defined in Lemma 4.16. 4.6 Let the affine functions fn : R≥0 → R for n ∈ N be given by fn (t) = γn t+ηn , t ≥ 0, where γn , ηn > 0. Show that the function ϕ(t) := inf fn (t), t ≥ 0, is continuous, monotonic, and concave.

n∈N

4.7 (a) Show that the set Uδ from Corollary 4.18 is open in L2 (0, 1). (b) Let K : Uδ → H 2 (0, 1) be the operator from Theorem 4.19; that is, c → u where u ∈ H 2 (0, 1) solves the boundary value problem (4.11). Show that this K satisfies the Tangential Cone Condition of Assumption 4.32. Hint: Modify the differential equation (4.15) such that u − u ˜ appears on the right hand side and continue as in the proof of Theorem 4.19.

Chapter 5

Inverse Eigenvalue Problems 5.1

Introduction

Inverse eigenvalue problems are not only interesting in their own right, but also have important practical applications. We recall the fundamental paper by Kac [148]. Other applications appear in parameter identification problems for parabolic or hyperbolic differential equations—as we study in Section 5.6 for a model problem—(see also [167, 187, 255]) or in grating theory ([156]). We study the Sturm–Liouville eigenvalue problem in canonical form. The direct problem is to determine the eigenvalues λ and the corresponding eigenfunctions u = 0 such that −

d2 u(x) + q(x) u(x) = λ u(x) , dx2

0 ≤ x ≤ 1,

u(0) = 0 and hu (1) + Hu(1) = 0 , 2

2

(5.1a) (5.1b)

2

where q ∈ L (0, 1) and h, H ∈ R with h +H > 0 are given. In this chapter, we assume that all functions are real-valued. In some applications, e.g., in grating theory, complex-valued functions q are also of practical importance. Essentially, all of the results of this chapter hold also for complex-valued q and are proven mainly by the same arguments. We refer to the remarks at the end of each section. The eigenvalue problem (5.1a), (5.1b) is a special case of the more general eigenvalue problem to determine ρ ∈ R and non-vanishing w such that     d dw(t) p(t) + ρ r(t) − g(t) w(t) = 0 , t ∈ [a, b] , (5.2a) dt dt αa w (a) + βa w(a) = 0 ,

αb w (b) + βb w(b) = 0 .

© Springer Nature Switzerland AG 2021 A. Kirsch, An Introduction to the Mathematical Theory of Inverse Problems, Applied Mathematical Sciences 120, https://doi.org/10.1007/978-3-030-63343-1 5

(5.2b)

169

170

Inverse Eigenvalue Problems

Here p, r, and g are given functions with p(t) > 0 and r(t) > 0 for t ∈ [a, b], and αa , αb , βa , βb ∈ R are constants with αa2 + βa2 > 0 and αb2 + βb2 > 0. If we assume, however, that g ∈ C[a, b] and p, r ∈ C 2 [a, b], then the Liouville transformation reduces the eigenvalue problem (5.2a), (5.2b) to the canonical form (5.1a), (5.1b). In particular, we define  σ(t) :=

r(t) , p(t)



1/4 f (t) := p(t) r(t) ,

b L :=

σ(s) ds ,

(5.3)

a

the monotonic function x : [a, b] → [0, 1] by x(t) :=

1 L

t σ(s) ds ,

t ∈ [a, b] ,

(5.4)

a

  and the new function u : [0, 1] → R by u(x) := f t(x) w t(x) , x ∈ [0, 1], where t = t(x) denotes the inverse of x = x(t). Elementary calculations show that u satisfies the differential equation (5.1a) with λ = L2 ρ and

   f (t) p(t)f  (t) 2 g(t) + . (5.5) q(x) = L r(t) r(t) f (t)2 t=t(x)

Also, it is easily checked that the boundary conditions (5.2b) are mapped into the boundary conditions h0 u (0) + H0 u(0) = 0 and h1 u (1) + H1 u(1) = 0 (5.6)  with h0 = αa σ(a)/ L f (a) and H0 = βa /f (a)−αa f  (a)/f (a)2 and, analogously, h1 , H1 with a replaced by b. In this chapter, we restrict ourselves to the study of the canonical Sturm– Liouville eigenvalue problem (5.1a), (5.1b). In the first part, we study the case h = 0 in some detail. At the end of Section 5.3, we briefly discuss the case where h = 1. In Section 5.3, we prove that there exists a countable number of eigenvalues λn of this problem and also prove an asymptotic formula. Because q is real-valued, the problem is self-adjoint, and the existence of a countable number of eigenvalues follows from the general spectral theorem of functional analysis (see Appendix A.6, Theorem A.53). Because this general theorem provides only the information that the eigenvalues tend to infinity, we need other tools to obtain more information about the rate of convergence. The basic ingredient in the proof of the asymptotic formula is the asymptotic behavior of the fundamental system of the differential equation (5.1a) as |λ| tends to infinity. Although all of the data and the eigenvalues are real-valued, we use results from complex analysis, in particular, Rouch´e’s theorem. This makes it necessary to allow the parameter λ in the fundamental system to be complex-valued. The existence of a fundamental solution and its asymptotics is the subject of the next section.

5.2

Construction of a Fundamental System

171

Section 5.5 is devoted to the corresponding inverse problem: Given the eigenvalues λn , determine the function q. In Section 5.6, we demonstrate how inverse spectral problems arise in a parameter identification problem for a parabolic initial value problem. Section 5.7, finally, studies numerical procedures for recovering q that have been suggested by Rundell and others (see [186, 233, 234]). We finish this section with a “negative” result, as seen in Example 5.1. Example 5.1 Let λ be an eigenvalue and u a corresponding eigenfunction of −u (x) + q(x) u(x) = λ u(x), 0 < x < 1,

u(0) = 0, u(1) = 0 .

Then λ is also an eigenvalue with corresponding eigenfunction v(x) := u(1 − x) of the eigenvalue problem −v  (x) + q˜(x) v(x) = λ v(x), 0 < x < 1,

v(0) = 0, v(1) = 0 ,

where q˜(x) := q(1 − x). This example shows that it is generally impossible to recover the function q unless more information is available. We will see that q can be recovered uniquely, provided we know that it is an even function with respect to 1/2 or if we know a second spectrum; that is, a spectrum for a boundary condition different from u(1) = 0.

5.2

Construction of a Fundamental System

It is well-known from the theory of linear ordinary differential equations that the following initial value problems are uniquely solvable for every fixed (realor complex-valued) q ∈ C[0, 1] and every given λ ∈ C: −u1 + q(x) u1 = λ u1 , 0 < x < 1 ,

u1 (0) = 1 , u1 (0) = 0

(5.7a)

−u2 + q(x) u2 = λ u2 , 0 < x < 1 ,

u2 (0) = 0 , u2 (0) = 1 .

(5.7b)

Uniqueness and existence for q ∈ L2 (0, 1) is shown in Theorem 5.4 below. The set of functions {u1 , u2 } is called a fundamental system of the differential equation −u + q u = λ u in (0, 1). The functions u1 and u2 are linearly independent because the Wronskian determinant is one

u1 u2 (5.8) [u1 , u2 ] := det = u1 u2 − u1 u2 = 1 . u1 u2 This is seen from d [u1 , u2 ] = u1 u2 − u1 u2 = u1 (q − λ) u2 − u2 (q − λ) u1 = 0 dx and [u1 , u2 ](0) = 1. The functions u1 and u2 depend on λ and q. We express this dependence often by uj = uj (·, λ, q), j = 1, 2. For q ∈ L2 (0, 1), the solution

172

Inverse Eigenvalue Problems

is not twice continuously differentiable anymore but is only an element of the Sobolev space   x 1  H (0, 1) := u ∈ C [0, 1] : u (x) = α + v(t) dt, α ∈ C, v ∈ L2 (0, 1) , 2

0

see (1.24). We write u for v and observe that u ∈ L2 (0, 1). The most important example is when q = 0. In this case, we can solve (5.7a) and (5.7b) explicitly and have the following: Example 5.2 Let q = 0. Then the solutions of (5.7a) and (5.7b) are given by √ sin( λ x) √ , u1 (x, λ, 0) = cos( λ x) and u2 (x, λ, 0) = λ √

(5.9)

respectively. An arbitrary branch of the square root can be taken because s → cos(sx) and s → sin(sx)/s are even functions. We will see that the fundamental solution for any function q ∈ L2 (0, 1) behaves as (5.9) as |λ| tends to infinity. For the proof of the next theorem, we need the following technical lemma. Lemma 5.3 Let q ∈ L2 (0, 1) and k, k˜ ∈ C[0, 1] such that there exists μ > 0 ˜ )| ≤ exp(μτ ) for all τ ∈ [0, 1]. Let K, K ˜ : with |k(τ )| ≤ exp(μτ ) and |k(τ C[0, 1] → C[0, 1] be the Volterra integral operators with kernels k(x − t) q(t) and ˜ − t) q(t), respectively; that is, k(x x k(x − t) q(t) x(t) dt ,

(Kφ)(x) =

0 ≤ x ≤ 1,

0

˜ Then the following estimate holds: and analogously for K.   ˜ K n−1 φ)(x) ≤ φ∞ 1 qˆ(x)n eμx , 0 ≤ x ≤ 1 , (K (5.10) n! x for all φ ∈ C[0, 1] and all n ∈ N. Here, qˆ(x) := 0 |q(t)| dt. If φ ∈ C[0, 1] satisfies also the estimate |φ(τ )| ≤ exp(μτ ) for all τ ∈ [0, 1], then we have   ˜ K n−1 φ)(x) ≤ 1 qˆ(x)n eμx , (K n!

0 ≤ x ≤ 1,

for all n ∈ N. Proof: We prove the estimates by induction with respect to n.

(5.11)

5.2

Construction of a Fundamental System

173

For n = 1, we estimate   ˜  (Kφ)(x)

x    ˜ − t) q(t) φ(t) dt =  k(x  0

≤

x

φ∞

eμ(x−t) |q(t)| dt ≤ φ∞ eμx qˆ(x) . 0

˜ Now we assume the validity of (5.10) for n. Because it holds also for K = K, we estimate x   ˜ K n φ)(x) ≤ (K eμ(x−t) |q(t)| |(K n φ)(t)| dt 0

≤

φ∞

1 μx e n!

We compute the last integral by x x n |q(t)| qˆ(t) dt = qˆ (t) qˆ(t)n dt = 0

0

=

x |q(t)| qˆ(t)n dt . 0

1 n+1

x 0

d qˆ(t)n+1 dt dt

1 qˆ(x)n+1 . n+1

This proves the estimate (5.10) for n + 1. For estimate (5.11), we only change the initial step n = 1 into   ˜  ≤ (Kφ)(x)

x eμ(x−t) eμt |q(t)| dt ≤ eμx qˆ(x) . 0

The remaining part is proven by the same arguments.



Now we prove the equivalence of the initial value problems for uj , j = 1, 2, to Volterra integral equations. Theorem 5.4 Let q ∈ L2 (0, 1) and λ ∈ C. Then we have (a) u1 , u2 ∈ H 2 (0, 1) are solutions of (5.7a) and (5.7b), respectively, if and only if u1 , u2 ∈ C[0, 1] solve the Volterra integral equations: √ x √ sin λ(x − t) √ u1 (x) = cos( λ x) + q(t) u1 (t) dt, (5.12a) λ 0

u2 (x)

=

√ √ x sin λ(x − t) sin( λ x) √ √ + q(t) u2 (t) dt , λ λ 0

respectively, for 0 ≤ x ≤ 1.

(5.12b)

174

Inverse Eigenvalue Problems

(b) The integral equations (5.12a) and (5.12b) and the initial value problems (5.7a) and (5.7b) are uniquely solvable. The solutions can be represented by a Neumann series. Let K denote the integral operator x (Kφ)(x) := 0

√ sin λ(x − t) √ q(t) φ(t) dt , λ

x ∈ [0, 1] ,

(5.13)

and define √ C(x) := cos( λx) Then u1 =

∞ 

Kn C

and

S(x) :=

and

u2 =

n=0

∞ 

√ sin( λx) √ . λ

(5.14)

Kn S .

(5.15)

n=0

The series converge uniformly with respect to (x, λ, q) ∈ [0, 1] × Λ × Q for all bounded sets Λ ⊂ C and Q ⊂ L2 (0, 1). Proof: (a) We use the following version of partial integration for f, g ∈ H 2 (0, 1): b

    b f (t) g(t) − f (t) g  (t) dt = f  (t) g(t) − f (t) g  (t) a .

(5.16)

a

We restrict ourselves to the proof for u1 . Let u1 be a solution of (5.7a). Then x

x S(x − t) q(t) u1 (t) dt =

0

x = 0

  S(x − t) λ u1 (t) + u1 (t) dt

0

  u1 (t) λ S(x − t) + S  (x − t) dt    =0

   t=x  + u1 (t) S(x − t) + u1 (t) S (x − t)  t=0 √ = u1 (x) − cos λx . On the other hand, let u1 ∈ C[0, 1] be a solution of the integral  xequation (5.12a). The operator A : L2 (0, 1) → L2 (0, 1) defined by (Aφ)(x) = 0 S(x − t) φ(t) dt, x ∈ (0, 1), is bounded, and it is easily seen that (Aφ) + λ(Aφ) = φ for φ ∈ C[0, 1]. Therefore, A is even bounded from L2 (0, 1) into H 2 (0, 1). Writing (5.12a) in the form u1 = C + A(qu1 ) yields that u1 ∈ H 2 (0, 1) and u1 = −λC + qu1 − λA(qu1 ) = qu1 − λu1 . This proves the assertion because the initial conditions are obviously satisfied. √ √ (b) We observe √ that all √ of the functions k(τ ) = cos( λτ ), k(τ ) = sin( λτ ), and k(τ ) = sin( λτ )/ λ for τ ∈ [0, 1] satisfy the estimate |k(τ )| ≤ exp(μ τ )

5.2

Construction of a Fundamental System

175

√ with μ = | Im λ|. This is obvious for the first two functions. For the third, it follows from √   τ τ √  sin( λτ )   √  ≤ | cos( λs)| ds = cosh(μs) ds ≤ cosh(μτ ) ≤ eμτ .  λ  0

0

We √ have to√study the integral operator K with kernel k(x−t)q(t), where k(τ ) = ˜ = K. Estimate (5.10) yields sin( λτ )/ λ. We apply Lemma 5.3 with K K n ∞ ≤

qˆ(1)n μ e < 1 n!

for sufficiently large n uniformly for q ∈ Q and λ ∈ Λ. Therefore, the Neumann series converges (see Appendix A.3, Theorem A.31), and part (b) is proven.  The integral representation of the previous theorem yields the following asymptotic behavior of the fundamental system by comparing the case for arbitrary q with the case of q = 0. Theorem 5.5 Let q ∈ L2 (0, 1), λ ∈ C, and u1 , u2 be the fundamental system; that is, the solutions of the initial value problems (5.7a) and (5.7b), respectively. Then we have for all x ∈ [0, 1]: √   u1 (x) − cos( λx) √     u2 (x) − sin(√ λx)   λ  √ √    u1 (x) + λ sin( λx)

≤

  x √ 1 √ exp | Im λ| x + |q(t)| dt , | λ|

(5.17a)

≤

  x √ 1 exp | Im λ| x + |q(t)| dt , |λ|

(5.17b)

0

0

≤

  x √ exp | Im λ| x + |q(t)| dt ,

(5.17c)

0

√    u2 (x) − cos( λx)

≤

  x √ 1 √ exp | Im λ| x + |q(t)| dt . | λ|

(5.17d)

0

√ Proof: Again, we √ use√the Neumann series and define C(τ ) := cos( λτ ) and Let K be the integral operator with kernel S(τ ) := √ λ. √ sin( λτ )/ q(t) sin λ(x − t) / λ. Then ∞  √ |(K n C)(x)| . |u1 (x) − cos( λx)| ≤ n=1

√ √ √ ˜ ) = sin( λτ ) and k(τ ) = sin( λτ )/ λ and denote by K ˜ and Now we set k(τ ˜ K the Volterra integral operators with kernels k(x−t) and k(x−t), respectively.

176

Inverse Eigenvalue Problems

Then K n =

√1 λ

˜ K n−1 and, by Lemma 5.3, part (b), we conclude that K

1 |(K C)(x)| ≤ √ | λ| n!

n

x |q(t)| dt

n

√  exp | Im λ| x

0

for n ≥ 1. Summation now yields the desired estimate: √ |u1 (x) − cos( λx)| ≤

  x √ 1 √ exp | Im λ| x + |q(t)| dt . | λ| 0

√  Because |S(x)| ≤ |√1λ| exp | Im λ| x , the same arguments prove the estimate (5.17b). Differentiation of the integral equations (5.12a) and (5.12b) yields u1 (x)

x √ √ √ + λ sin( λ x) = cos λ(x − t) q(t) u1 (t) dt, 0

√ u2 (x) − cos( λ x) =

x cos

√ λ(x − t) q(t) u2 (t) dt .

0

˜ defined as the operator with kernel q(t) cos With K as before and K Then u1 (x) +

√ λ(x − t).

∞  √ √ ˜ λ sin( λ x) = K K nC , n=0

∞  √ ˜ u2 (x) − cos( λ x) = K K nS , n=0

and we use Lemma 5.3, estimate (5.11), again. Summation yields the estimates (5.17c) and (5.17d).  In the next section, we need the fact that the eigenfunctions are continuously differentiable with respect to q and λ. We remind the reader of the concept of Fr´echet differentiability (F-differentiability) of an operator between Banach spaces X and Y (see Appendix A.7, Definition A.60). Here we consider the mapping (λ, q) → uj (·, λ, q) from C × L2 (0, 1) into C[0, 1] for j = 1, 2. We denote these mappings by uj again and prove the following theorem: Theorem 5.6 Let uj : C × L2 (0, 1) → C[0, 1], j = 1, 2, be the solution operator of (5.7a) and (5.7b), respectively. Then we have the following: (a) uj is continuous.

5.2

Construction of a Fundamental System

177

ˆ qˆ) ∈ C × L2 (0, 1) with (b) uj is continuously F-differentiable for every (λ, partial derivatives ∂ ˆ qˆ) = uj,λ (·, λ, ˆ qˆ) uj (·, λ, (5.18a) ∂λ and ∂ ˆ qˆ) (q) = uj,q (·, λ, ˆ qˆ) , uj (·, λ, (5.18b) ∂q ˆ qˆ) and uj,q (·, λ, ˆ qˆ) are solutions of the following initial where uj,λ (·, λ, boundary value problems for j = 1, 2:  ˆ uj,λ = uj (·, λ, ˆ qˆ) in (0, 1) −uj,λ + qˆ − λ uj,λ (0) = 0, uj,λ (0) = 0,  ˆ uj,q = −q uj (·, λ, ˆ qˆ) + qˆ − λ

−uj,q

uj,q (0) = 0,

(5.19) in (0, 1),

uj,q (0) = 0.

(c) Furthermore, for all x ∈ [0, 1] we have x

uj (t)2 dt

=

[uj,λ , uj ](x),

u1 (t) u2 (t) dt

=

[u1,λ , u2 ](x) = [u2,λ , u1 ](x),

(5.20b)

q(t) uj (t)2 dt

=

[uj,q , uj ](x),

(5.20c)

q(t) u1 (t) u2 (t) dt

=

[u1,q , u2 ](x) = [u2,q , u1 ](x) ,

j = 1, 2,

(5.20a)

0

x 0

x

− x −

j = 1, 2,

0

(5.20d)

0

where [u, v] denotes the Wronskian determinant from (5.8). Proof: (a), (b): Continuity and differentiability of uj follow from the integral equations (5.12a) and (5.12b) because the kernel and the right-hand sides depend continuously and differentiably on λ and q. It remains to show the representation of the derivatives in (b). Let u = uj , j = 1 or 2. Then ˆ + ε) −u (·, λ  ˆ −u (·, λ)

+ +



ˆ − ε u(·, λ ˆ + ε) qˆ − λ  ˆ u(·, λ) ˆ qˆ − λ

= 0, = 0;

thus −

    1 ˆ + ε) − u(·, λ) ˆ  + qˆ − λ ˆ + ε) − u(·, λ) ˆ = u(·, λ ˆ + ε) . ˆ 1 u(·, λ u(·, λ ε ε

178

Inverse Eigenvalue Problems

Furthermore, the homogeneous initial conditions are satisfied for the difference ˆ as ε → 0. Therequotient. The right-hand side converges uniformly to u(·, λ) fore, the difference quotient converges to uλ uniformly in x. The same arguments yield the result for the derivative with respect to q. (c) Multiplication of the differential equation for uj,λ by uj and the differential equation for uj by uj,λ and subtraction yields u2j (x)

= uj (x) uj,λ (x) − uj,λ (x) uj (x) =

d   u (x) uj,λ (x) − uj,λ (x) uj (x) . dx j

Integration of this equation and the homogeneous boundary conditions yield the first equation of (5.20a). The proofs for the remaining equations use the same arguments and are left to the reader.  At no place in this section have we used the assumption that q is real-valued. Therefore, the assertions of Theorems 5.4, 5.5, and 5.6 also hold for complexvalued q.

5.3

Asymptotics of the Eigenvalues and Eigenfunctions

We first restrict ourselves to the Dirichlet problem; that is, the eigenvalue problem −u (x) + q(x) u(x) = λ u(x) , 0 < x < 1,

u(0) = u(1) = 0 .

(5.21)

We refer to the end of this section for different boundary conditions. Again, let q ∈ L2 (0, 1) be real-valued. We observe that λ ∈ C is an eigenvalue of this problem if and only if λ is a zero of the function f (λ) := u2 (1, λ, q) . Again, u2 = u2 (·, λ, q) denotes the solution of the differential equation −u2 + q u2 = λ u2 in (0, 1) with initial conditions u2 (0) = 0 and u2 (0) = 1. If u2 (1, λ, q) = 0, then u = u2 (·, λ, q) is an eigenfunction corresponding to the eigenvalue λ, normalized such that u (0) = 1. There are different ways to normalize the eigenfunctions. Later we will sometimes normalize them such that the L2 −norms are one; that is, use g = u/uL2 instead of u. The function f plays exactly the role of the well-known characteristic polynomial for matrices and is, therefore, called the characteristic function of the eigenvalue problem. Theorem 5.6 implies that f is differentiable; that is, analytic in all of C. This observation makes it possible to use tools from complex analysis. First, we summarize well-known facts about eigenvalues and eigenfunctions for the Sturm–Liouville problem. Theorem 5.7 Let q ∈ L2 (0, 1) be real-valued. Then

5.3

Asymptotics of the Eigenvalues and Eigenfunctions

179

(a) All eigenvalues λ are real. (b) There exists a countable number of real eigenvalues λj , j ∈ N, which tend to infinity as j → ∞. The corresponding eigenfunctions gj ∈ C[0, 1], normalized by gj L2 = 1, form a complete orthonormal system in L2 (0, 1). (c) The geometric and algebraic multiplicities of the eigenvalues λj are one; that is, the eigenspaces are one-dimensional and the zeros of the characteristic function f are simple. (d) Let the eigenvalues be ordered as λ1 < λ2 < λ3 < · · · . The eigenfunction gj corresponding to λj has exactly j − 1 zeros in (0, 1). (e) Let q be even with respect to 1/2; that is, q(1 − x) = q(x) for all x ∈ [0, 1]. Then gj is even with respect to 1/2 for odd j and odd with respect to 1/2 for even j. Proof: (a) and (b) follow from the fact that the boundary value problem is self-adjoint. We refer to Problem 5.1 for a repetition of the proof (see also Theorems A.52 and A.53). (c) Let λ be an eigenvalue and u, v be two corresponding eigenfunctions. Choose α, β with α2 + β 2 > 0, such that α u (0) = β v  (0). The function w := αu − βv solves the differential equation and w(0) = w (0) = 0; that is, w vanishes identically. Therefore, u and v are linearly dependent. We apply Theorem 5.6, part (c), to show that λ is a simple zero of f . Because u2 (1, λ, q) = 0, we have from (5.20a) for j = 2 that f  (λ)

= =

∂ u2 (1, λ, q) = u2,λ (1, λ, q) ∂λ 1 1 u2 (x, λ, q)2 dx = 0 . u2 (1, λ, q)

(5.22)

0

This proves part (c). (d) First we note that every gj has only a finite number of zeros. Otherwise, they would accumulate at some point x ∈ [0, 1], and it is not difficult to show that gj and also gj vanish at x. This would imply that gj vanishes identically. We fix j ∈ N and define the function h : [0, 1] × [0, 1] → R by h(t, x) = uj (x; tq). Here, uj (·; tq) is the j-th eigenfunction uj corresponding to tq instead of q and normalized such that uj (0) = 1. Then h is continuously differentiable and h(t, 0) = h(t, 1) = 0 and every zero of h(t, ·) is simple. This holds for every t. By part (a) of Lemma 5.8 below, the number of zeros of h(t, ·) is constant with respect to t. Therefore, √ uj (·, q) = h(1, ·) has exactly the same number of zeros as h(0, x) = uj (x, 0) = 2 sin(jπx) which is j − 1. The normalization to gj = uj (·, q)/uj (·, q)L2 does not change the number of zeros. (e) Again, it is sufficient to prove that vj = u2 (·, λj , q) is even (or odd) for odd (or even) j. First, we note that also v˜j (x) := vj (1 − x) is an eigenfunction.

180

Inverse Eigenvalue Problems

Since the eigenspace is one-dimensional, we conclude that there exists ρj ∈ R with v˜j (x) = vj (1 − x) = ρj vj (x) for all x. For x = 1/2 this implies that (1 − ρj )vj (1/2) = 0 and also by differentiation (1 + ρj )vj (1/2) = 0 for all j. Since vj (1/2) and vj (1/2) cannot vanish simultaneously, we conclude that ρj ∈ {+1, −1} and even ρj = ρj vj (0) = −vj (1). From (5.22), we conclude that sign f  (λj ) = sign vj (1) = − sign ρj . Since λj are the subsequent zeros of f , we conclude that f  (λj )f  (λj+1 ) < 0 for all j; that is, ρj = σ(−1)j+1 , j ∈ N, for some σ ∈ {+1, −1}. The first eigenfunction v1 has no zero by part (d) which yields σ = 1. This ends the proof.  The first part of the following technical result has been used in the previous proof, the second part will be needed below. Lemma 5.8 (a) Let h : [0, 1] × [0, 1] → R be continuously differentiable such that h(t, ·) has finitely many zeros in [0, 1] and all are simple for every t ∈ [0, 1]. Then the number m(t) of zeros of h(t, ·) is constant with respect to t. (b) Let z ∈ C with |z − nπ| ≥ π/4 for all n ∈ Z. Then  exp | Im z| < 4 | sin z| . Proof: (a) It suffices to show that t → m(t) is continuous. Fix tˆ ∈ [0, 1] and ˆj )/∂x = 0 there let x ˆj , j = 1, . . . , m(tˆ), be the zeros of h(tˆ, ·). Because ∂h(tˆ, x xj − δ, x ˆj + δ) ∩ [0, 1] with exist intervals T = (tˆ − δ, tˆ + δ) ∩ [0, 1]and Jj = (ˆ ∂h(t, x)/∂x = 0 for all t ∈ T and x ∈ j Jj . Therefore, for every t ∈ T , the function h(t, ·) has at most one zero in every Jj . On the other hand, by the implicit function theorem (applicable because ∂h(tˆ, x ˆj )/∂x = 0), for every t ∈ T , where T and Jj are possibly there exists at least one zero of h(t, ·) in every J j  made smaller. Outside of Jj there are no zeros of h(tˆ, ·), and thus by making perhaps T smaller again, no zeros of h(t, ·) either for all t ∈ T . This shows that m(t) = m(tˆ) for all t ∈ T which shows that t → m(t) is continuous, thus constant. / {nπ : (b) Let ψ(z) = exp |z2 |/| sin z| for z = z1 + iz2 , z1 , z2 ∈ R with z1 ∈ n ∈ Z}. We consider two cases: 1st case: |z2 | > ln 2/2. Then ψ(z) =

2 e|z2 | 2 e|z2 | 2 ≤ = < 4 iz −z −iz +z |z | −|z | 1 2 1 2 2 2 |e −e | e −e 1 − e−2|z2 |

because exp(−2|z2 |) < 1/2. 2nd case: |z2 | ≤ ln 2/2. From |z − nπ| ≥ π/4 for all n, we conclude that |z1 − nπ|2 ≥ π 2 /16 − z22 ≥ π 2 /16 − (ln 2)2 /4 ≥ π 2 /64; thus | sin z1 | ≥ sin π8 . With | Re sin z| = | sin z1 | | cosh z2 | ≥ | sin z1 |, we conclude that √ √ 2 2 e|z2 | ≤ ≤ < 4.  ψ(z) ≤ | Re sin z| | sin z1 | | sin π8 |

5.3

Asymptotics of the Eigenvalues and Eigenfunctions

181

Now we prove the “counting lemma,” a first crude asymptotic formula for the eigenvalues. As the essential tool in the proof, we use the theorem of Rouch´e from complex analysis (see [2]), which we state for the convenience of the reader: Let U ⊂ C be a domain and the functions F and G be analytic in C and |F (z) − G(z)| < |G(z)| for all z ∈ ∂U . Then F and G have the same number of zeros in U.  Lemma 5.9 Let q ∈ L2 (0, 1) and N > 2 exp qL1 be an integer. Then (a) The characteristic function f (λ) := u2 (1, λ, q) has exactly N zeros in the half-plane   (5.23) H := λ ∈ C : Re λ < (N + 1/2)2 π 2 . (b) For every m > N there exists exactly one zero of f in the set √    Um := λ ∈ C :  λ − mπ  < π/2 . √ Here we take the branch with Re λ ≥ 0.

(5.24)

(c) There are no other zeros of f in C. Proof: We are going to apply Rouch´e’s theorem to the function F (z) = f (z 2 ) = u2 (1, z 2 , q) and the corresponding function G of the eigenvalue problem for q = 0; that is, G(z) := sin z/z. For U , we take one of the sets Wm or VR defined by   Wm := z ∈ C : |z − mπ| < π/2 ,   z ∈ C : | Re z| < (N + 1/2)π, | Im z| < R VR := for fixed R > (N + 1/2)π and want to apply Lemma 5.8: (i) First let z ∈ ∂Wm : For n ∈ Z, n = m, we have |z−nπ| ≥ |m−n|π−|z−mπ| ≥ π − π/2 > π/4. For n = m, we observe that |z − mπ| = π/2 > π/4. Therefore, we can apply Lemma 5.8 for z ∈ ∂Wm . Furthermore, we note the estimate |z| ≥ mπ − |z − mπ| = (m − 1/2)π > N π > 2N for all z ∈ ∂Wm . (ii) Let z ∈ ∂VR , n ∈ Z. Then | Re z| = (N +1/2)π or | Im z| = R. In either case, we estimate |z − nπ|2 = (Re z − nπ)2 + (Im z)2 ≥ π 2 /4 > π 2 /16. Therefore, we can apply Lemma 5.8 for z ∈ ∂VR . Furthermore, we have the estimate |z| ≥ (N + 1/2)π > 2N for all z ∈ ∂VR . Application of Theorem 5.5 and Lemma 5.8 yields the following estimate for all z ∈ ∂VR ∪ ∂Wm :      1 4 | sin z| N F (z) − sin z  ≤ exp | Im z| + qL1 ≤  z  |z|2 |z|2 2      sin z  2N  sin z  . <  = |z|  z  z  Therefore, F and G(z) := sin z/z have the same number of zeros in VR and every Wm . Because the zeros of G are ±nπ, n = 1, 2, . . ., we conclude that G

182

Inverse Eigenvalue Problems

has exactly 2N zeros in VR and exactly one zero in every Wm . By the theorem of Rouch´e, this also holds for F .  Now we show that F has no  zero outside of VR ∪ m>N Wm . Again, we apply Lemma 5.8: Let z ∈ / VR ∪ m>N Wm . From z ∈ / VR , we conclude that  |z| = (Re z)2 + (Im z)2 ≥ (N +1/2)π. For n > N , we have that |z −nπ| > π/2 because z ∈ / Wn . For n ≤ N , we conclude that |z − nπ| ≥ |z| − nπ ≥ (N + 1/2 − n)π ≥ π/2. We apply Theorem 5.5 and Lemma 5.8 again and use the second triangle inequality. This yields    sin z    − 1 exp | Im z| + qL1  |F (z)| ≥   2 z |z|   

  sin z  1 4 exp q L  1− ≥  z  |z|  

 sin z   1 − 2N > 0 ≥  z  |z| because |z| ≥ (N + 1/2)π > 2N . Therefore, we have shown that f has exactly one zero in every Um , m > N , and N zeros in the set √  √    HR := λ ∈ C : 0 < Re λ < (N + 1/2)π, Im λ < R and no other zeros. It√remains  to show that HR ⊂ H. For λ = |λ| exp(iθ)∈ H R, we conclude that Re λ = |λ| cos θ2 < (N + 1/2)π; thus Re λ = |λ| cos 2 θ2 ≤  |λ| cos2 θ2 < (N + 1/2)2 π 2 . This lemma proves again the existence of infinitely many eigenvalues. The arguments are not changed for the case of complex-valued functions q. In this case, the general spectral theory is not applicable anymore because the boundary value problem is not self-adjoint. This lemma also provides more information about the eigenvalue distribution, even for the real-valued case. First, we order the eigenvalues in the form λ 1 < λ2 < λ3 < · · · . Lemma 5.9 implies that  λn = nπ + O(1) ;

that is, λn = n2 π 2 + O(n) .

(5.25)

For the treatment of the inverse problem, it is necessary to improve this formula. It is our aim to prove that 2 2

λn = n π

1 ˜n q(t) dt + λ

+ 0

where

∞   2 ˜ λn  < ∞ .

(5.26)

n=1

There are several methods to prove (5.26). We follow the treatment in [218]. The key is to apply the fundamental theorem of calculus to the function t → λn (tq)

5.3

Asymptotics of the Eigenvalues and Eigenfunctions

183

for t ∈ [0, 1], thus connecting the eigenvalues λn corresponding to q with the eigenvalues n2 π 2 corresponding to q = 0 by the parameter t. For this approach, we need the differentiability of the eigenvalues with respect to q. For fixed n ∈ N, the function q → λn (q) from L2 (0, 1) into C is well-defined and Fr´echet differentiable by the following theorem. Theorem 5.10 For every n ∈ N, the mapping q → λn (q) from L2 (0, 1) into C is continuously Fr´echet differentiable for every qˆ ∈ L2 (0, 1) and λn (ˆ q )q

1 =

gn (x, qˆ)2 q(x) dx ,

q ∈ L2 (0, 1) .

(5.27)

0

Here, ˆ n , qˆ) u2 (x, λ  gn (x, qˆ) :=  ˆ n , qˆ) 2 u2 (·, λ L ˆ n := λn (ˆ denotes the L2 -normalized eigenfunction corresponding to λ q ). Note 2 ˆ that the integral is well-defined because u2 (·, λn , qˆ) ∈ H (0, 1) ⊂ C[0, 1]. ˆ n , qˆ) = 0 and apply the implicit function theorem Proof: We observe that u2 (1, λ to the equation u2 (1, λ, q) = 0 ˆ n , qˆ). This is possible because the zero λ ˆ n of u2 (1, ·, qˆ) in a neighborhood of (λ is simple by Lemma 5.7. The implicit function theorem (see Appendix A.7, Theorem A.66) yields the existence of a unique function λn = λn (q) such that u2 (1, λn (q), q) = 0 for all q in a neighborhood of qˆ; we know this already. But it also implies that the function λn is continuously differentiable with respect to q and ∂ ∂ ˆ n , qˆ) λ (ˆ ˆ n , qˆ)q ; 0 = u2 (1, λ u2 (1, λ n q )q + ∂λ ∂q that is, u2,λ (1) λn (ˆ q )q + u2,q (1) = 0. With Theorem 5.6, part (c), we conclude that λn (ˆ q )q

u2,q (1) u2 (1) u2,q (1) = − u2,λ (1) u2,λ (1) u2 (1) 1 1 q(x) u2 (x)2 dx [u2,q , u2 ](1) 0 = = gn (x, qˆ)2 q(x) dx , = − 1 [u2,λ , u2 ](1) u2 (x)2 dx

= −

0

ˆ and qˆ. where we have dropped the arguments λ

0



Now we are ready to formulate and prove the main theorem which follows:

184

Inverse Eigenvalue Problems

Theorem 5.11 Let Q ⊂ L2 (0, 1) be bounded, q ∈ Q, and λn ∈ C the corresponding eigenvalues. Then we have 2 2

λn = n π

1

1 q(t) dt −

+ 0

q(t) cos(2nπt) dt + O(1/n)

(5.28)

0

for n → ∞ uniformly for q ∈ Q. Furthermore, the corresponding eigenfunctions gn , normalized to gn L2 = 1, have the following asymptotic behavior: √ gn (x) = 2 sin(nπx) + O(1/n) and (5.29a) √  gn (x) = 2 n π cos(nπx) + O(1) (5.29b) as n → ∞ uniformly for x ∈ [0, 1] and q ∈ Q. We observe that the second integral on the right-hand side of (5.28) isthe nth Fourier coefficient an of q with respect to cos(2πnt) : n = 0, 1, 2, . . . , . From Fourier theory, it is known that ∞an converges to zero, and even more: Bessel’s inequality (A.7) yields that n=0 |an |2 < ∞; that is, (5.26) is satisfied. If q is smooth enough, e.g., continuously differentiable, then an tends to zero faster than 1/n. In that case, this term is absorbed in the O(1/n) expression. Proof: We split the proof into four parts: √ √ (a) First, we show that gn (x) = 2 sin( λn x) √ + O(1/n) uniformly for (x, q) ∈ [0, 1] × Q. By Lemma 5.9, we know that λn = nπ + O(1), and thus by Theorem 5.5 √ sin( λn x) √ + O(1/n2 ) . u2 (x, λn ) = λn 1 With the formula 2 0 sin2 (αt) dt = 1 − sin(2α)/(2α), we compute 1

2

u2 (t, λn ) dt = 0

= =

1 λn

1

sin2



λn t dt + O(1/n3 )

0 √

1 sin(2 λn ) √ 1− + O(1/n3 ) 2λn 2 λn 1 [1 + O(1/n)] . 2λn

Therefore, we have  √ u2 (x, λn ) = 2 sin( λn x) + O(1/n) . gn (x) =  1 u (t, λn )2 dt 0 2 (b) Now we show that

√

λn = nπ + O(1/n) and gn (x) =

√ 2 sin(nπx) + O(1/n).

5.3

Asymptotics of the Eigenvalues and Eigenfunctions

185

We apply the fundamental theorem of calculus and use Theorem 5.10 2 2

λn − n π

1 = λn (q) − λn (0) = 1 =

λn (tq)q dt =

0

d λn (tq) dt dt

0 1 1

0

(5.30)

gn (x, tq)2 q(x) dx dt = O(1) .

0

√ This √ yields λn = nπ+O(1/n) and, with part (a), the asymptotic form gn (x) = 2 sin(nπx) + O(1/n). (c) Now the asymptotics of the eigenvalues follow easily from (5.30) by the observation that gn (x, tq)2 = 2 sin2 (nπx) + O(1/n) = 1 − cos(2nπx) + O(1/n) , uniformly for t ∈ [0, 1] and q ∈ Q. (d) Similarly, we have for the derivatives √ √ √ 2 λn cos( λn x) + O(1) u (x, λn )  gn (x) =  2 = 1 1 + O(1/n) u (t, λn )2 dt 0 2 =

√

2 nπ cos(nπx) + O(1) .



Example 5.12 We illustrate Theorem 5.11 by the following two numerical examples:  (a) Let q1 (x) = exp sin(2πx) , x ∈ [0, 1]. Then q1 is analytic and periodic with period 1. Plots√of the √ characteristic functions λ → f (λ) for q1 and q = 0; that is, λ → sin λ/ λ are shown in Figure 5.1. 1.4

0.5

1.2

0.4

1

0.3

0.8

0.2

0.6 0.1 0.4 0

0.2

-0.1

0

-0.2

-0.2 -0.4 0

2

4

6

8

10

12

14

16

18

20

-0.3 0

10

20

30

40

50

60

70

80

90

100

Figure 5.1: Characteristic functions of q, q1 , respectively, on [0, 20] and [5, 100]

186

Inverse Eigenvalue Problems

(b) Let q2 (x) = −5 x for 0 ≤ x ≤ 0.4 and q2 (x) = 4 for 0.4 < x ≤ 1. The function q2 is not continuous. Plots of the characteristic functions λ → f (λ) for q2 and q = 0 are shown in Figure 5.2. 1.4

0.6

1.2

0.5

1

0.4

0.8

0.3

0.6

0.2

0.4

0.1

0.2

0

0

-0.1

-0.2

-0.2

-0.4 0

2

4

6

8

10

12

14

16

18

20

-0.3 0

10

20

30

40

50

60

70

80

90

100

Figure 5.2: Characteristic functions of q, q2 , respectively, on [0, 20] and [5, 100] The Fourier coefficients of q1 converge to zero of exponential order. The following table shows the eigenvalues λn corresponding to q1 , the eigenvalues n2 π 2 corresponding to q = 0 and the difference 1 2 2 q(x) dx for n = 1, . . . , 10 : cn := λn − n π − 0

λn 11.1 40.9 90.1 159.2 248.0 356.6 484.9 632.9 800.7 988.2

2 2

n π 9.9 39.5 88.8 157.9 246.7 354.3 483.6 631.7 799.4 987.0

cn −2.04 ∗ 10−2 1.49 ∗ 10−1 2.73 ∗ 10−3 −1.91 ∗ 10−3 7.74 ∗ 10−4 4.58 ∗ 10−4 4.58 ∗ 10−4 4.07 ∗ 10−4 3.90 ∗ 10−4 3.83 ∗ 10−4

We clearly observe the rapid convergence. Because q2 is not continuous, the Fourier coefficients converge to zero only slowly. Again, we list the eigenvalues λn for q2 , the eigenvalues n2 π 2 corresponding to q = 0, and the differences 1 2 2 q(x) dx and cn := λn − n π − 0

dn

:=

2 2

λn − n π

1 −

1 q(x) dx +

0

q(x) cos(2πnx) dx 0

5.3

Asymptotics of the Eigenvalues and Eigenfunctions

187

for n = 1, . . . , 10: λn 12.1 41.1 91.1 159.8 248.8 357.4 484.5 633.8 801.4 989.0

n2 π 2 9.9 39.5 88.8 157.9 246.7 354.3 483.6 631.7 799.4 987.0

cn 1.86 ∗ 10−1 −3.87 ∗ 10−1 3.14 ∗ 10−1 1.61 ∗ 10−1 2.07 ∗ 10−2 8.29 ∗ 10−2 −1.25 ∗ 10−1 1.16 ∗ 10−1 −6.66 ∗ 10−2 5.43 ∗ 10−3

dn −1.46 ∗ 10−1 8.86 ∗ 10−2 2.13 ∗ 10−2 −6.70 ∗ 10−3 2.07 ∗ 10−2 −4.24 ∗ 10−3 6.17 ∗ 10−3 3.91 ∗ 10−3 −1.38 ∗ 10−3 5.43 ∗ 10−3

Now we sketch the modifications necessary for Sturm–Liouville eigenvalue problems of the type −u (x) + q(x) u(x) = λ u(x) , u(0) = 0 ,

0 < x < 1,

u (1) + Hu(1) = 0 .

(5.31a) (5.31b)

Now the eigenvalues are zeros of the characteristic function f (λ) = u2 (1, λ, q) + Hu2 (1, λ, q),

λ ∈ C.

(5.32) √ √ For the special case, where q = 0, we have u2 (x, λ, 0) = sin( λx)/ λ. The characteristic function for this case is then given by √ √ sin λ √ . g(λ) = cos λ + H λ The zeros of f for q = 0 and H = 0 are λn = (n + 1/2)2 π 2 , n = 0, 1, 2, . . . If H = 0, one has to solve the transcendental equation z cot z + H = 0. One can show (see Problem 5.2) by an application of the implicit function theorem in R2 that the eigenvalues for q = 0 behave as λn = (n + 1/2)2 π 2 + 2H + O(1/n) . Lemma 5.7 is also valid because the boundary value problem is again selfadjoint. The Counting Lemma 5.9 now takes the following form:  Lemma 5.13 Let q ∈ L2 (0, 1) and N > 2 exp qL1 (1 + |H|) be an integer. Then we have (a) The mapping f (λ) := u2 (1, λ, q) + H u2 (1, λ, q) has exactly N zeros in the half-plane   H := λ ∈ C : Re λ < N 2 π 2 .

188

Inverse Eigenvalue Problems

(b) f has exactly one zero in every set √    Um := λ ∈ C :  λ − (m − 1/2)π  < π/2 provided m > N . (c) There are no other zeros of f in C. For the proof, we refer to Problem 5.3. We can apply the implicit function theorem to the equation u2 (1, λn (q), q) + Hu2 (1, λn (q), q) = 0 because the zeros are again simple. Differentiating this equation with respect to q yields   ˆ n , qˆ) + Hu2,λ (1, λ ˆ n , qˆ) λ (ˆ u2,λ (1, λ n q )q ˆ n , qˆ) + Hu2,q (1, λ ˆ n , qˆ) = 0 . + u2,q (1, λ Theorem 5.6 yields 1

u2 (t)2 dt

u2 (1) − u2,λ (1) u2 (1)   

= u2,λ (1)

0

=−Hu2 (1)

  = −u2 (1) u2,λ (1) + H u2,λ (1) ˆ n and qˆ. Analogously, we compute where again we have dropped the arguments λ 1 −

  q(t)u2 (t)2 dt = −u2 (1) u2,q (1) + H u2,q (1)

0

and thus q )q λn (ˆ

u2,q (1) + H u2,q (1) = = −  u2,λ (1) + H u2,λ (1)

1

q(t)u2 (t)2 dt . 1 u (t)2 dt 0 2

0

This has the same form as before. We continue as in the case of the Dirichlet boundary condition and arrive at Theorem 5.14. Theorem 5.14 Let Q ⊂ L2 (0, 1) be bounded, q ∈ Q, and H ∈ R. The eigenvalues λn have the asymptotic form  λn =

1 n+ 2

2

2

1

1 q(t) dt −

π + 2H + 0

q(t) cos(2n + 1)πt dt + O(1/n) (5.33) 0

5.4

Some Hyperbolic Problems

189

as n tends to infinity, uniformly in q ∈ Q. For the L2 -normalized eigenfunctions, we have √ 2 sin(n + 1/2)πx + O(1/n) and (5.34a) gn (x) = √  2 (n + 1/2) π cos(n + 1/2)πx + O(1) (5.34b) gn (x) = uniformly for x ∈ [0, 1] and q ∈ Q. As mentioned at the beginning of this section, there are other ways to prove the asymptotic formulas for the eigenvalues and eigenfunctions that avoid Lemma 5.9 and the differentiability of λn with respect to q. But the proof in, e.g., [276], seems to yield only the asymptotic behavior λn = m2n π 2 +

1 q(t) dt + O(1/n) 0

instead of (5.28). Here, (mn ) denotes some sequence of natural numbers. Before we turn to the inverse problem, we make some remarks concerning the case where q is complex-valued. Now the eigenvalue problems are no longer self-adjoint, and the general spectral theory is not applicable anymore. With respect to Lemma 5.7, it is still easy to show that the eigenfunctions corresponding to different eigenvalues are linearly independent and that the geometric multiplicities are still one. The Counting Lemma 5.9 is valid without restrictions. From this, we observe also that the algebraic multiplicities of λn are one, at least for n > N . Thus, the remaining arguments of this section are valid if we restrict ourselves to the eigenvalues λn with n > N . Therefore, the asymptotic formulas (5.28), (5.29a), (5.29b), (5.33), (5.34a), and (5.34b) hold equally well for complex-valued q.

5.4

Some Hyperbolic Problems

As a preparation for the following sections, in particular, Sections 5.5 and 5.7, we study some initial value problems for the two-dimensional linear hyperbolic partial differential equation ∂ 2 W (x, t) ∂ 2 W (x, t) − + a(x, t) W (x, t) = 0 , ∂x2 ∂t2 where the coefficient a has the special form a(x, t) = p(t)−q(x). It is well-known that the method of characteristics reduces initial value problems to Volterra integral equations of the second kind, which can be studied in spaces of continuous functions. This approach naturally leads to solution concepts for nonsmooth coefficients and boundary data. We summarize the results in three theorems. In each of them, we formulate first the results for the case of smooth coefficients and then for the nonsmooth case. We remark that it is not our aim to

190

Inverse Eigenvalue Problems

relax the solution concept to the weakest possible case but rather to relax the assumptions only to the extent that are needed in Sections 5.5 and 5.7 and in Subsection 7.6.3. Although most of the problems—at least for smooth data—are subjects of elementary courses on partial differential equations, we include the complete proofs for the convenience of the reader. Before we begin with the statements of the theorems, we recall some function spaces   C00 [0, 1] := f ∈ C[0, 1] : f (0) = 0 ,    x 1 2 H (0, 1) := f ∈ C[0, 1] : f (x) = α + g(t) dt, α ∈ R, g ∈ L (0, 1) , 1 (0, 1) H00

0

:= H 1 (0, 1) ∩ C00 [0, 1]

and equip them with their canonical norms f ∞

:=

f H 1

:=

max |f (x)| in C00 [0, 1],  1 f 2L2 + f  2L2 in H 1 (0, 1) and H00 (0, 1) .

0≤x≤1

1 The notations C00 [0, 1] and H00 (0, 1) should indicate that the boundary condition is set only at x = 0. By  · C j for j ≥ 1 we denote the canonical norm in C j [0, 1]. Furthermore, we define the triangular regions Δ0 ⊂ R2 and Δ ⊂ R2 by   Δ0 := (x, t) ∈ R2 : 0 < t < x < 1 , (5.35a)   2 Δ := (x, t) ∈ R : |t| < x < 1 , (5.35b)

respectively. We begin with an initial value problem, sometimes called the Goursat problem. Theorem 5.15 (a) Let p, q ∈ C[0, 1] and f ∈ C 2 [0, 1] with f (0) = 0. Then there exists a unique solution W ∈ C 2 (Δ0 ) of the following hyperbolic initial value problem:  ∂ 2 W (x, t) ∂ 2 W (x, t) − + p(t) − q(x) W (x, t) = 0 ∂x2 ∂t2 W (x, x)

= f (x) ,

W (x, 0)

=

0,

0 ≤ x ≤ 1,

0 ≤ x ≤ 1.

in Δ0 ,

(5.36a) (5.36b) (5.36c)

(b) The solution operator (p, q, f ) → W has an extension to a bounded operator from L2 (0, 1) × L2 (0, 1) × C00 [0, 1] into C(Δ0 ).  (c) The operator (p, q, f ) → W (1, ·), Wx (1, ·) has an extension to a bounded 1 (0, 1) into H 1 (0, 1) × L2 (0, 1). Here and operator from L2 (0, 1) × L2 (0, 1) × H00 in the following, we denote by Wx the partial derivative with respect to x.

5.4

Some Hyperbolic Problems

191

Proof: (a) First, we extend the problem to the larger region Δ and study the problem ∂ 2 W (x, t) ∂ 2 W (x, t) − + a(x, t) W (x, t) = 0 in Δ , ∂x2 ∂t2 W (x, x) W (x, −x)

= f (x) , 0 ≤ x ≤ 1 , = −f (x) , 0 ≤ x ≤ 1 ,

(5.37a) (5.37b) (5.37c)

where we have extended p(t) − q(x) to a(x, t) := p(|t|) − q(x) for (x, t) ∈ Δ. To treat problem (5.37a)–(5.37c), we make the change of variables x = ξ + η,

t = ξ −η.

Then (x, t) ∈ Δ if and only if (ξ, η) ∈ D, where   D := (ξ, η) ∈ (0, 1) × (0, 1) : η + ξ < 1 .

(5.38)

We set w(ξ, η) := W (ξ+η, ξ−η) for (ξ, η) ∈ D. Then W solves problem (5.37a)– (5.37c) if and only if w solves the hyperbolic problem ∂ 2 w(ξ, η) = −a(ξ + η, ξ − η) w(ξ, η) ,    ∂ξ ∂η

(ξ, η) ∈ D ,

(5.39a)

=:˜ a(ξ,η)

w(ξ, 0) w(0, η)

= f (ξ) for ξ ∈ [0, 1] , = −f (η) for η ∈ [0, 1] .

(5.39b) (5.39c)

Now let w be a solution of (5.39a)–(5.39c). We integrate the differential equation twice and use the initial conditions. Then w solves the integral equation η ξ w(ξ, η) = 0

a ˜(ξ  , η  ) w(ξ  , η  ) dξ  dη  − f (η) + f (ξ) ,

(5.40)

0

for (ξ, η) ∈ D. This is a Volterra integral equation in two dimensions. We use the standard method to solve this equation by successive iteration in C(D) ˜ ∈ L2 (D). Let A be the Volterra where we assume only p, q ∈ L2 (0, 1), and thus a integral operator defined by the integral on the right-hand side of (5.40). By induction with respect to n ∈ N, it can easily be seen (compare with the proof of Theorem 5.18 below for a similar, but more complicated estimate) that  1 (ξ η)n/2 , | An w (ξ, η)| ≤ w∞ ˜ anL2 n!

n = 1, 2, . . . ;

1 thus An w∞ ≤ w∞ ˜ anL2 n! . Therefore, An L(C[0,1]) < 1 for sufficiently large n, and the Neumann series converges (see Appendix A.3, Theorem A.31).

192

Inverse Eigenvalue Problems

This proves that there exists a unique solution w ∈ C(D) of (5.40). From our arguments, uniqueness also holds for (5.37a)–(5.37c). Now we prove that the solution w ∈ C(D) is even in C 2 (D). Obviously, from (5.40) and the differentiability of f , we conclude that w is differentiable with partial derivative (with respect to ξ, the derivative with respect to η is seen analogously) η wξ (ξ, η)

=

  q(ξ + η  ) − p(|ξ − η  |) w(ξ, η  ) dη  + f  (ξ)

0

=

ξ+η  ξ q(s) w(ξ, s − ξ) ds − p(|s|) w(ξ, ξ − s) ds + f  (ξ) . ξ

ξ−η

This second form can be differentiated again. Thus w ∈ C 2 (D), and we have shown that W is the unique solution of (5.37a)–(5.37c). Because a(x, ·) is an even function and the initial data are odd functions with respect to t, we conclude from the uniqueness result that the solution W (x, ·) is also odd. In particular, this implies that W (x, 0) = 0 for all x ∈ [0, 1], which proves that W solves problem (5.36a)–(5.36c) and finishes part (a). Part (b) follows immediately from the integral equation (5.40) because the ˜ ∈ integral operator A : C(D) → C(D) depends continuously on the kernel a L2 (D). For part (c), we observe that W (1, 2ξ − 1) d W (1, 2ξ − 1) dξ

= w(ξ, 1 − ξ) , thus d = w(ξ, 1 − ξ) = wξ (ξ, 1 − ξ) − wη (ξ, 1 − ξ) dξ 1 1 wξ (ξ, 1 − ξ) + wη (ξ, 1 − ξ) . = 2 2

Wx (1, 2ξ − 1)

and

We have computed wξ already above and have 1 wξ (ξ, 1 − ξ) =

ξ q(s) w(ξ, s − ξ) ds −

ξ

p(|s|) w(ξ, ξ − s) ds + f  (ξ)

2ξ−1

1 (0, 1). An analogous formula which is in L2 (0, 1) for p, q ∈ L2 (0, 1) and f ∈ H00 1 holds for wη and shows that W (t, ·) ∈ H (0, 1) and Wx (1, ·) ∈ L2 (0, 1) for 1 (0, 1). This ends the proof.  p, q ∈ L2 (0, 1) and f ∈ H00

Remark 5.16 (a) If p, q ∈ L2 (0, 1) and f ∈ C[0, 1] with f (0) = 0, we call the solution W ∈ C(Δ0 ), given by   1 1 W (x, t) = w (x + t), (x − t) , 2 2

5.4

Some Hyperbolic Problems

193

where w ∈ C(D) solves the integral equation (5.40), the weak solution of the Goursat problem (5.36a)–(5.36c). We observe that for every weak solution W there exist sequences (pn ), (qn ) in C[0, 1] and (fn ) in C 2 [0, 1] with fn (0) = 0 and pn − pL2 → 0, qn − qL2 → 0 and fn − f ∞ → 0 such that the solutions Wn ∈ C 2 (Δ0 ) of (5.36a)–(5.36c) corresponding to pn , qn , and fn converge uniformly to W . (b) We observe from the integral equation (5.40) that w has a decomposition into w(ξ, η) = w1 (ξ, η) − f (η) + f (ξ) where w1 ∈ C 1 (D) even if only p, q ∈  L2 (0, 1). This transforms into W (x, t) = W1 (x, t) − f 12 (x − t) + f 12 (x + t) with W1 ∈ C 1 (Δ). For the special case p = q = 0, the integral equation (5.40) reduces to the well-known solution formula ! ! 1 1 W (x, t) = f (x + t) − f (x − t) . 2 2 The next theorem studies a Cauchy problem for the same hyperbolic differential equation. Theorem 5.17 (a) Let f ∈ C 2 [0, 1], g ∈ C 1 [0, 1] with f (0) = f  (0) = g(0) = 0, and p, q ∈ C[0, 1] and F ∈ C 1 (Δ0 ). Then there exists a unique solution W ∈ C 2 (Δ0 ) of the Cauchy problem  ∂ 2 W (x, t) ∂ 2 W (x, t) − + p(t) − q(x) W (x, t) = F (x, t) 2 2 ∂x ∂t W (1, t) ∂ W (1, t) ∂x

in Δ0 , (5.41a)

= f (t)

for 0 ≤ t ≤ 1,

(5.41b)

= g(t)

for 0 ≤ t ≤ 1 .

(5.41c)

(b) Furthermore, the solution operator (p, q, F, f, g) → W has an extension to a 1 bounded operator from L2 (0, 1) × L2 (0, 1) × L2 (Δ0 ) × H00 (0, 1) × L2 (0, 1) into C(Δ0 ). Proof: As in the proof of Theorem 5.15, we set a(x, t) := p(|t|) − q(x) for (x, t) ∈ Δ and extend F to an even function on Δ by F (x, −t) = F (x, t) for (x, t) ∈ Δ0 . We also extend f and g to odd functions  on [−1, 1] by f (−t) = −f (t) and g(−t) = −g(t) for t ∈ [0, 1]. Then F ∈ C 1 Δ \ ([0, 1] × {0}) ∩ C(Δ), f ∈ C 2 [−1, 1], and g ∈ C 1 [−1, 1]. We again make the change of variables x = ξ + η,

t = ξ − η,

w(ξ, η) = W (ξ + η, ξ − η) for (ξ, η) ∈ D ,

where D is given by (5.38). Then W solves (5.41a)–(5.41c) if and only if w solves ∂ 2 w(ξ, η) = a ˜(ξ, η) w(ξ, η) + F˜ (ξ, η), (ξ, η) ∈ D , ∂ξ ∂η

194

Inverse Eigenvalue Problems

where F˜ (ξ, η) = F (ξ+η, ξ−η) and a ˜(ξ, η) = −a(ξ+η, ξ−η) = q(ξ+η)−p(|ξ−η|. The Cauchy conditions (5.41b) and (5.41c) transform into w(ξ, 1 − ξ) = f (2ξ − 1)

and wξ (ξ, 1 − ξ) + wη (ξ, 1 − ξ) = 2 g(2ξ − 1)

for 0 ≤ ξ ≤ 1. Differentiating the first equation and solving for wξ and wη yields wξ (ξ, 1 − ξ) = g(2ξ − 1) + f  (2ξ − 1)

and wη (ξ, 1 − ξ) = g(2ξ − 1) − f  (2ξ − 1)

for 0 ≤ ξ ≤ 1. Integration of the differential equation with respect to ξ from ξ to 1 − η yields ∂w(ξ, η) = − ∂η

1−η 

   a ˜(ξ , η) w(ξ  , η) + F˜ (ξ  , η) dξ  + g(1 − 2η) − f  (1 − 2η) .

ξ

Now we integrate this equation with respect to η from η to 1 − ξ and arrive at 1−ξ 1−η  

w(ξ, η)

= η





 a ˜(ξ  , η  ) w(ξ  , η  ) + F˜ (ξ  , η  ) dξ  dη 

(5.42)

ξ

1−ξ  1 1 − g(1 − 2η  ) dη  + f (2ξ − 1) + f (1 − 2η) 2 2 η

for (ξ, η) ∈ D. This is again a Volterra integral equation in two variables. The solution w ∈ C(D) is in C 2 (D). Indeed, since it is obviously differentiable we take the derivative with respect to η and arrive at the formula above and, after substitution of a ˜(ξ  , η) and making the change of variables s = ξ  − η and s = ξ  + η, respectively, at the representation 1−2η  1   ∂w(ξ, η) = q(s) w(s − η, η) ds p(|s|) w(s + η, η) ds − ∂η ξ−η

ξ+η 1−η 

F˜ (ξ  , η) dξ  + g(1 − 2η) − f  (1 − 2η) .

− ξ

An analogous formula holds for the derivative with respect to ξ. We can differ 1−η F˜ (ξ  , η) dξ  is differentiable in entiate again because the function ψ(ξ, η) = ξ D – although F˜ is not differentiable at the line ξ = η. The reader should try to 1 (0, 1), and g ∈ L2 (0, 1) prove this. If only p, q ∈ L2 (0, 1), F ∈ L2 (Δ0 ), f ∈ H00 then (5.42) defines the weak solution. Let A denote the integral operator 1−ξ 1−η  

(Aw)(ξ, η) = η

ξ



a ˜(ξ  , η  ) w(ξ  , η  ) dξ  dη  ,

(ξ, η) ∈ D .

5.4

Some Hyperbolic Problems

195

By induction, it is easily seen that |(An w)(ξ, η)| ≤ w∞ ˜ anL2 



1 (2n)!

1−ξ−η

n

for all (ξ, η) ∈ D and n ∈ N; thus anL2  An w∞ ≤ w∞ ˜

1 (2n)!

for all n ∈ N. For sufficiently large n, we conclude that An L(C[0,1]) < 1, which again implies that (5.42) is uniquely solvable in C(D) for any p, q, g ∈ L2 (0, 1), 1 (0, 1).  F ∈ L2 (Δ0 ), and f ∈ H00 For the special case p = q = 0 and F = 0, the integral equation (5.42) reduces to the well-known d’Alembert formula 1 W (x, t) = − 2

t+(1−x) 

g(τ ) dτ +

1  1  f t + (1 − x) + f t − (1 − x) . 2 2

t−(1−x)

Finally, the third theorem studies a quite unusual coupled system for a pair (W, r) of functions. We treat this system with the same methods as above. Theorem 5.18 (a) Let q ∈ C[0, 1], F ∈ C 1 (Δ0 ), f ∈ C 2 [0, 1], and g ∈ C 1 [0, 1] such that f (0) = f  (0) = g(0) = 0. Then there exists a unique pair of functions (W, r) ∈ C 2 (Δ0 ) × C 1 [0, 1] with ∂ 2 W (x, t) ∂ 2 W (x, t) − − q(x) W (x, t) = F (x, t) r(x) 2 ∂x ∂t2 W (x, x) W (x, 0)

= =

1 2

in Δ0 ,

(5.43a)

x

0,

r(s) ds ,

0 ≤ x ≤ 1,

(5.43b)

0

0 ≤ x ≤ 1,

(5.43c)

and W (1, t) = f (t)

and

∂ W (1, t) = g(t) ∂x

for all t ∈ [0, 1] .

(5.43d)

(b) Furthermore, the solution operator (q, F, f, g) → (W, r) has an extension to 1 a bounded operator from L2 (0, 1) × C(Δ0 ) × H00 (0, 1) × L2 [0, 1] into C(Δ0 ) × 2 L (0, 1).

196

Inverse Eigenvalue Problems

Proof: (a) We apply the same arguments as in the proofs of Theorems 5.15 and 5.17. We extend F (x, ·) to an even function and f and g to odd functions. We again make the change of variables x = ξ + η and t = ξ − η and set F˜ (ξ, η) = F (ξ + η, ξ − η). In Theorem 5.17 (for p = 0), we have shown that the solution W of the Cauchy problem (5.43a) and (5.43d) is equivalent to the integral equation 1−ξ 1−η  

w(ξ, η)



= η

   q(ξ + η  ) w(ξ  , η  ) + F˜ (ξ  , η  ) r(ξ  + η  ) dξ  dη 

ξ 1−ξ 

g(1 − 2η  ) dη  +

−

1 1 f (2ξ − 1) + f (1 − 2η) 2 2

(5.44a)

η

for w(ξ, η) = W (ξ + η, ξ − η) (see equation (5.42)). From this and the initial condition (5.43b), we derive a second integral equation. We set η = 0 in (5.44a), differentiate, and substitute (5.43b). This yields the following Volterra equation after an obvious change of variables: 1 r(x) 2

1 = −

  q(s) w(x, s − x) + r(s) F˜ (x, s − x) ds

x

+ g(2x − 1) + f  (2x − 1) .

(5.44b) 2

Assume that there exists a solution (w, r) ∈ C(D) × L (0, 1) of (5.44a) and (5.44b). From (5.44b), we observe that r is continuous on [0, 1], and thus by (5.44a), w ∈ C 1 (D) and thus also r ∈ C 1 [0, 1] because the function ψ(x) = 1 r(s) F˜ (x, s − x) ds is differentiable on [0, 1]. Therefore, the right-hand side x of (5.43a) is differentiable and we conclude as in the previous theorem that d d W (x, x) = dx w(x, 0) = 12 r(x). Now, because w ∈ C 2 (D). Furthermore, dx F (x, ·) is even and f and g are odd functions, we conclude that W (x, ·) is also odd. In particular,W (x, 0) = 0 for all x ∈ [0, 1]. This implies W (0, 0) = 0 and x thus W (x, x) = 12 0 r(s) ds. Therefore, we have shown that every solution of equations (5.44a) and (5.44b) satisfies (5.43a)–(5.43d) and vice versa. Now we sketch the proof that the system (5.44a), (5.44b) is uniquely solvable 1 (0, 1), for (w, r) ∈ C(D) × L2 (0, 1) for given q ∈ L2 (0, 1), F ∈ C(Δ0 ), f ∈ H00 2 and g ∈ L (0, 1). This would also include the proof of part (b). We write this system in the form (w, r) = A(w, r) + b in the product space C(D) × L2 (0, 1) 2 which we equip with the norm (w, r)∞,L2 = max{w∞ , r  L (0,1) }. To apply the fixed point theorem we define first the constant c := 2 qL2 (0,1) + F ∞ and, for given (w, r) ∈ C(D) × L2 (0, 1), the functions (wn , rn ) = An (w, r). By induction we prove the following estimates n   wn (ξ, η) ≤ (w, r)∞,L2 √c (1 − ξ − η)n/2 , (ξ, η) ∈ D , n! n   c rn (x) ≤ (w, r)∞,L2 √ (1 − x)n/2 , x ∈ (0, 1) , n!

5.5

The Inverse Problem

197

for all n = 1, 2, . . ..  1−ξ  1−η We use the elementary integral η (1 − ξ  − η  )n dξ  dη  = ξ 1 n+2 for n = 0, 1, . . . and set q˜(ξ, η) = q(ξ + η) and r˜(ξ, η) (n+1)(n+2) (1 − ξ − η)  = r(ξ + η) for abbreviation. We note that ˜ q 2L2 (D) = D |q(ξ + η)|2 d(ξ, η) = 1 t |q(t)|2 dt ≤ q2L2 (0,1) and analogously ˜ rL2 (D) ≤ rL2 (D) . 0 For n = 1 we have by the Cauchy-Schwarz inequality   w1 (ξ, η)

1−ξ 1−η  

≤

w∞



1−ξ 1−η  



|˜ q (ξ, η)| dξ dη + F ∞ η

≤



η

ξ

≤ ≤ ≤ ≤

|˜ r(ξ, η)| dξ  dη 

ξ

" # 1−ξ 1−η    # # q L2 (D) + F ∞ $ (w, r)∞,L2 ˜ dξ  dη  η

  r(x)



ξ

 1 c (w, r)∞,L2 √ (1 − ξ − η) ≤ c (w, r)∞,L2 1 − ξ − η , 2  √ 2 w∞ qL2 (0,1) + rL2 (0,1) F ∞ 1 − x  √ 2 (w, r)∞,L2 qL2 (0,1) + F ∞ 1 − x √ c (w, r)∞,L2 1 − x .

The step from n to n+1 is proven in just the same way. Therefore, An ∞,L2 ≤ cn √ which tends to zero as n tends to infinity. Application of Theorem A.31 n!

yields that (5.44a), (5.44b) has a unique solution in C(D) × L2 (0, 1) for all 1 (0, 1), and g ∈ L2 (0, 1).  q ∈ L2 (0, 1), F ∈ C(Δ0 ), f ∈ H00

5.5

The Inverse Problem

Now we study the inverse spectral problem. This is, given the eigenvalues λn of the Sturm–Liouville eigenvalue problem −u (x) + q(x) u(x) = λ u(x) , 0 < x < 1 ,

u(0) = 0 , u(1) = 0 ,

(5.45)

determine the function q. We saw in Example 5.1 that the knowledge of the spectrum {λn : n ∈ N} is, in general, not sufficient to determine q uniquely. We need more information, such as a second spectrum μn of an eigenvalue problem of the form −v  (x) + q(x) v(x) = μ v(x) ,

v(0) = 0 , v  (1) + Hv(1) = 0 ,

(5.46)

or some knowledge about the eigenfunctions. The basic tool in the uniqueness proof for this inverse problem is the use of the Gelfand–Levitan–Marchenko integral operator (see [101]). This integral operator maps solutions of initial value problems for the equation −u +qu = λu onto solutions for the equation −u + pu = λu and, most importantly, does not

198

Inverse Eigenvalue Problems

depend on λ. It turns out that the kernel of this operator is the solution for the hyperbolic boundary value problem that was studied in the previous section. Theorem 5.19 Let p, q ∈ L2 (0, 1), λ ∈ C, and u, v ∈ H 2 (0, 1) be solutions of −u (x) + q(x) u(x) = λ u(x) ,

0 < x < 1,

u(0) = 0 ,

(5.47a)

−v  (x) + p(x) v(x) = λ v(x) ,

0 < x < 1,

v(0) = 0 ,

(5.47b)

such that u (0) = v  (0). Also let K ∈ C(Δ0 ) be the weak solution of the Goursat problem  ∂ 2 K(x, t) ∂ 2 K(x, t) − + p(t) − q(x) K(x, t) = 0 2 2 ∂x ∂t K(x, 0) K(x, x)

=

0,

=

1 2

in Δ0 ,

0 ≤ x ≤ 1, x



(5.48a) (5.48b)

q(s) − p(s) ds ,

0 ≤ x ≤ 1,

(5.48c)

0

where the triangular region Δ0 is again defined by   Δ0 := (x, t) ∈ R2 : 0 < t < x < 1 .

(5.49)

Then we have x u(x) = v(x) +

K(x, t) v(t) dt ,

0 ≤ x ≤ 1.

(5.50)

0

 x We remark that Theorem 5.15 with f (x) = 12 0 q(s) − p(s) ds implies that this Goursat problem is uniquely solvable in the weak sense. Proof: First, let p, q ∈ C[0, 1]. Then K ∈ C 2 (Δ0 ) by Theorem 5.15. Define w by the right-hand side of (5.50); that is, x K(x, t) v(t) dt

w(x) := v(x) +

for 0 ≤ x ≤ 1 .

0

Then w(0) = v(0) = 0 = u(0) and w is differentiable with 



x

w (x) = v (x) + K(x, x)v(x) +

Kx (x, t) v(t) dt, 0

0 < x < 1.

5.5

The Inverse Problem

199

Again, we denote by Kx , Kt , etc., the partial derivatives. For x = 0, we have w (0) = v  (0) = u (0). Furthermore, d K(x, x) + K(x, x) v  (x) dx x + Kx (x, x) v(x) + Kxx (x, t) v(t) dt

w (x) = v  (x) + v(x)

0



d = p(x) − λ + K(x, x) + Kx (x, x) v(x) + K(x, x) v  (x) dx x   (q(x) − p(t))K(x, t)v(t) + Ktt (x, t) v(t) dt . + 0

Partial integration yields x Ktt (x, t) v(t) dt 0

x

 t=x K(x, t) v  (t) dt + Kt (x, t) v(t) − K(x, t) v  (t) t=0

= 0

x =



p(t) − λ K(x, t) v(t) dt + Kt (x, x) v(x) − K(x, x) v  (x) .

0

Therefore, we have w (x)

=



d K(x, x) + Kx (x, x) + Kt (x, x) v(x) p(x) − λ + dx   d =2 dx K(x,x)=q(x)−p(x)



+ q(x) − λ

=







⎡

x K(x, t) v(t) dt 0

q(x) − λ ⎣v(x) +

x

⎤ K(x, t)v(t) dt⎦ =



q(x) − λ w(x) ;

0

that is, w solves the same initial value problem as u. The Picard–Lindel¨ of uniqueness theorem for initial boundary value problems yields w = u. Thus, we have proven the theorem for smooth functions p and q. Now let p, q ∈ L2 (0, 1). Then we choose functions (pn ), (qn ) in C[0, 1] with pn → p and qn → q in L2 (0, 1), respectively. Let Kn be the solution of (5.48a)–

200

Inverse Eigenvalue Problems

(5.48c) for pn and qn . We have already shown that x un (x) = vn (x) +

Kn (x, t) vn (t) dt,

0 ≤ x ≤ 1,

0

for all n ∈ N, where un and vn solve (5.47a) and (5.47b), respectively, with un (0) = vn (0) = u (0) = v  (0). From the continuous dependence results of Theorems 5.6 and 5.15(b), the functions un , vn , and Kn converge uniformly to u, v, and K, respectively. This proves the assertion of the theorem for  p, q ∈ L2 (0, 1). √ √ As an example, we take p = 0 and v(x) = sin( λx)/ λ and have the following result: Example 5.20 Let u be a solution of −u (x) + q(x) u(x) = λ u(x) ,

u (0) = 1 ,

u(0) = 0 ,

(5.51)

for given q ∈ L2 (0, 1). Then we have the representation u(x) =

√ √ x sin λ x sin λ t √ dt , + K(x, t) √ λ λ

0 ≤ x ≤ 1,

(5.52)

0

where the kernel K solves the following Goursat problem in the weak sense: Kxx (x, t) − Ktt (x, t) − q(x) K(x, t) = 0 K(x, 0) K(x, x)

=

0,

=

1 2

in Δ0 ,

0 ≤ x ≤ 1,

(5.53a) (5.53b)

x q(s) ds ,

0 ≤ x ≤ 1.

(5.53c)

0

This example has an application that is interesting in itself but that we also need in Section 5.7 and later in Subsection 7.6.3 Theorem 5.21 Let λn be the eigenvalues of one of the eigenvalue  problems √ (5.45) or (5.46) where again q ∈ L2 (0, 1). Then the set of functions sin λn · :   √ 1 n ∈ N is complete in L2 (0, 1). This means that 0 h(x) sin λn x dx = 0 for all n ∈ N implies that h = 0. Proof: Let T : L2 (0, 1) → L2 (0, 1) be the Volterra integral operator of the second kind with kernel K; that is, x (T v)(x) := v(x) +

K(x, t) v(t) dt, 0

x ∈ (0, 1), v ∈ L2 (0, 1) ,

5.5

The Inverse Problem

201

where K solves the Goursat problem (5.53a)–(5.53c) in the weak sense. Then we √ know that T is an isomorphism from L2 (0, 1) onto itself. Define vn (x) := sin λn x for x ∈ [0, 1], n ∈ N. Let un be the eigenfunction corresponding to λn , normalized to un (0) = 1. By the preceding example, 1 u n = √ T vn λn Now, if

1 0

or

vn =



λn T −1 un .

h(x) vn (x) dx = 0 for all n ∈ N, then 1

0 =

h(x) T −1 un (x) dx =

0

1

un (x) (T ∗ )−1 h(x) dx

for all n ∈ N ,

0

where T ∗ denotes the L2 -adjoint of T . Because {un /un L2 : n ∈ N} is complete  in L2 (0, 1) by Lemma 5.7, we conclude that (T ∗ )−1 h = 0 and thus h = 0. Now we can prove the main uniqueness theorem. Theorem 5.22 Let H ∈ R, p, q ∈ L2 (0, 1), and λn (p), λn (q) be the eigenvalues of the eigenvalue problem −u + r u = λ u in (0, 1),

u(0) = 0, u(1) = 0 ,

corresponding to r = p and r = q, respectively. Furthermore, let μn (p) and μn (q) be the eigenvalues of −u + r u = μ u in (0, 1),

u(0) = 0, u (1) + Hu(1) = 0 ,

corresponding to r = p and r = q, respectively. If λn (p) = λn (q) and μn (p) = μn (q) for all n ∈ N, then p = q. Proof: From the asymptotics of the eigenvalues (Theorem 5.11), we conclude that λn (p)

2 2

= n π

1 +

p(t) dt + o(1),

n → ∞,

q(t) dt + o(1) ,

n → ∞,

0

λn (q)

= n2 π 2 +

1 0

and thus

1



p(t) − q(t) dt = lim (λn (p) − λn (q)) = 0 . n→∞

(5.54)

0

Now let K be the weak solution of the Goursat problem (5.48a)–(5.48c). Then K depends only on p and q and is independent of the eigenvalues λn := λn (p) =

202

Inverse Eigenvalue Problems

λn (q) and μn := μn (p) = μn (q). Furthermore, from (5.54), we conclude that K(1, 1) = 0. Now let un , vn be the eigenfunctions corresponding to λn (q) and λn (p), respectively; that is, solutions of the differential equations −un (x) + q(x) un (x) = λn un (x),

−vn (x) + p(x) vn (x) = λn vn (x)

for 0 < x < 1 with homogeneous Dirichlet boundary conditions on both sides. Furthermore, we assume that they are normalized by un (0) = vn (0) = 1. Then Theorem 5.19 is applicable and yields the relationship x K(x, t) vn (t) dt for x ∈ [0, 1] , (5.55) un (x) = vn (x) + 0

and all n ∈ N. For x = 1, the boundary conditions yield 1 0 = K(1, t) vn (t) dt for all n ∈ N .

(5.56)

0

Now we use the fact that the set {vn /vn L2 : n ∈ N} forms a complete orthonormal system in L2 (0, 1). From this, K(1, t) = 0 for all t ∈ [0, 1] follows. Now let u ˜n and v˜n be eigenfunctions corresponding to μn and q and p, respectively, with the normalization u ˜n (0) = v˜n (0) = 1. Again, Theorem 5.19 is applicable and yields the relationship (5.55) for u ˜n and v˜n instead of un und vn , respectively. Assume for the moment that K is differentiable. Then we can differentiate this equation, set x = 1, and arrive at   0 = u ˜n (1) − v˜n (1) + H u ˜n (1) − v˜n (1) 1   Kx (1, t) + H K(1, t) v˜n (t) dt . = K(1, 1) v˜n (1) +       0

=0

1

=0

We conclude that 0 Kx (1, t) v˜n (t) dt = 0 for all n ∈ N. From this, Kx (1, t) = 0 vn L2 } forms a complete orthonormal for all t ∈ (0, 1) follows because {˜ vn /˜ system. However, K is only a weak solution since p, q ∈ L2 (0, 1). From part  1 (b) (x, t) − f of Remark 5.16 we know that K has the form K(x, t) = K 1 2 (x −  t) + f 12 (x + t) with K ∈ C 1 (Δ0 ) where in the present situation f (x) =  1 x 1 (approximation 2 0 (q(s) − p(s)) ds is in H00 (0, 1). Then one can easilyprove x  1 of f by C −functions) that ψ ∈ C 1 [0, 1] where ψ(x) = 0 f 12 (x ± t) v˜n (t) dt    x and ψ  (x) = f 12 (x ± x) v˜n (x) + 12 0 f  12 (x ± t) v˜n (t) dt. Therefore, one can argue as in the smooth case of K. Now we apply the uniqueness part of Theorem 5.17 (in particular, the integral equation (5.42) for f = g = 0 and F = 0) which yields that K has to vanish identically. In particular, this means that x 1  p(s) − q(s) ds for all x ∈ (0, 1) . 0 = K(x, x) = 2 0

5.6

A Parameter Identification Problem

203

Differentiating this equation yields that p = q.



We have seen in Example 5.1 that the knowledge of one spectrum for the Sturm–Liouville differential equation is not enough information to recover the function q uniquely. Instead of knowing the spectrum for a second pair of boundary conditions, we can use other kinds of information, as the following theorem shows: Theorem 5.23 Let p, q ∈ L2 (0, 1) with eigenvalues λn (p), λn (q), and eigenfunctions un and vn , respectively, corresponding to Dirichlet boundary conditions u(0) = 0, u(1) = 0. Let the eigenvalues coincide; that is, λn (p) = λn (q) for all n ∈ N. Let one of the following assumptions also be satisfied: (a) Let p and q be even functions with respect to 1/2; that is, p(1 − x) = p(x) and q(1 − x) = q(x) for all x ∈ [0, 1]. (b) Let the Neumann boundary values coincide; that is, let v  (1) un (1) = n  un (0) vn (0)

for all n ∈ N .

(5.57)

Then p = q. Proof: (a) From Theorem 5.7, part (e), we know that the eigenfunctions un and vn , again normalized by un (0) = vn (0) = 1, are even with respect to x = 1/2 for odd n and odd for even n. In particular, un (1) = vn (1). This reduces the uniqueness question for part (a) to part (b). (b) We follow the first part of the proof of Theorem 5.22. From (5.56), we again conclude that K(1, t) vanishes for all t ∈ (0, 1). The additional assumption (5.57) yields that un (1) = vn (1). We differentiate (5.55), set x = 1, 1 and arrive at 0 Kx (1, t)vn (t) dt = 0 for all n ∈ N. Again, this implies that Kx (1, ·) = 0, and the proof follows the same lines as the proof of Theorem 5.22. 

5.6

A Parameter Identification Problem

This section and the next two chapters are devoted to the important field of parameter identification problems for partial differential equations. In Chapter 6, we study the problem of impedance tomography to determine the conductivity distribution from boundary measurements, while in Chapter 7, we study the inverse scattering problem to determine the refractive index of a medium from measurements of the scattered field. In the present section, we consider an application of the inverse Sturm–Liouville eigenvalue problem to the following parabolic initial boundary value problem. First, we formulate the direct problem: Let T > 0 and ΩT := (0, 1) × (0, T ) ⊂ R2 , q ∈ C[0, 1] and f ∈ C 2[0, T ] be given with f (0) = 0 and q(x) ≥ 0 for x ∈ [0, 1]. Determine U ∈ C ΩT ,

204

Inverse Eigenvalue Problems

which is twice continuously differentiable with respect to x and continuously differentiable with respect to t in ΩT such that ∂U/∂x ∈ C ΩT and ∂U (x, t) ∂t

=

∂ 2 U (x, t) − q(x) U (x, t) ∂x2

U (x, 0)

=

0,

U (0, t)

=

0,

x ∈ [0, 1] , ∂ U (1, t) = f (t) , ∂x

in ΩT ,

(5.58a) (5.58b)

t ∈ (0, T ) .

(5.58c)

From the theory of parabolic initial boundary value problems, it is known that there exists a unique solution of this problem. We prove uniqueness and refer to [170] or (5.60) for the question of existence. Theorem 5.24 Let f = 0. Then U = 0 is the only solution of (5.58a)–(5.58c) in ΩT . Proof: Multiply the differential equation (5.58a) by U (x, t) and integrate with respect to x. This yields 1 d 2 dt

1

U (x, t)2 dx =

0

1 0

∂ 2 U (x, t) 2 U (x, t) − q(x) U (x, t) dx . ∂x2

We integrate by parts and use the homogeneous boundary conditions: 1 d 2 dt

1

1 

2

U (x, t) dx = − 0

0

∂U (x, t) ∂x

2

+ q(x) U (x, t) dx ≤ 0 . 2

1

This implies that t → 0 U (x, t)2 dx is nonnegative and monotonically nonin1 1 creasing. From 0 U (x, 0)2 dx = 0, we conclude that 0 U (x, t)2 dx = 0 for all t; that is, U = 0.  Now we turn to the inverse problem. Let f be known and, in addition, U (1, t) for all 0 < t ≤ T . The inverse problem is to determine the coefficient q. In this section, we are only interested in the question if this provides sufficient information in principle to recover q uniquely; that is, we study the question of uniqueness of the inverse problem. It is our aim to prove the following theorem: Theorem 5.25 Let U1 , U2 be solutions of (5.58a)–(5.58c) corresponding to q = q1 ≥ 0 and q = q2 ≥ 0, respectively, and to the same f ∈ C 2 [0, T ] with f (0) = 0 and f  (0) = 0. Let U1 (1, t) = U2 (1, t) for all t ∈ (0, T ). Then q1 = q2 on [0, 1]. Proof: Let (q, U ) be (q1 , U1 ) or (q2 , U2 ), respectively. Let λn and gn , n ∈ N, be the eigenvalues and eigenfunctions, respectively, of the Sturm–Liouville eigenvalue problem (5.46) for H = 0; that is, −u (x) + q(x) u(x) = λ u(x),

0 < x < 1,

u(0) = u (1) = 0 .

5.6

A Parameter Identification Problem

205

We assume that the eigenfunctions are normalized by gn L2 = 1 for all n ∈ N. Furthermore, we can assume that gn (1) > 0 for all n ∈ N1 . We know that {gn : n ∈ N} forms a complete orthonormal system in L2 (0, 1). Theorem 5.14 implies the asymptotic behavior λn

˜n (n + 1/2)2 π 2 + qˆ + λ

=

gn (x)

√  2 sin(n + 1/2)πx + O 1/n ,

=

∞ 

with

˜2 < ∞ , λ n

(5.59a)

n=1

(5.59b)

1

where qˆ = 0 q(x) dx. In the first step, we derive a series expansion for the solution U of the initial boundary value problem (5.58a)–(5.58c). From the completeness of {gn : n ∈ N}, we have the Fourier expansion U (x, t) =

∞ 

1 an (t) gn (x) with

U (x, t) gn (x) dx, n ∈ N ,

an (t) =

n=1

0

where the convergence is understood in the L2 (0, 1)-sense for every t ∈ (0, T ]. We would like to substitute this into the differential equation and the initial and boundary conditions. Because for this formal procedure the interchanging of summation and differentiation is not justified, we suggest a different derivation of an . We differentiate an and use the partial differential equation (5.58a). This yields an (t)

1 = 0

∂U (x, t) gn (x) dx = ∂t

1 0

∂ 2 U (x, t) − q(x) U (x, t) gn (x) dx ∂x2

x=1 ∂U (x, t) − U (x, t)gn (x) = gn (x) ∂x x=0 1    + U (x, t) gn (x) − q(x) gn (x) dx    0

=−λn gn (x)

= f (t) gn (1) − λn an (t) . With the initial condition an (0) = 0, the solution is given by t an (t) = gn (1)

f (τ ) e−λn (t−τ ) dτ ;

0

that is, the solution U of (5.58a)–(5.58c) takes the form U (x, t) =

∞ 

t gn (1) gn (x)

n=1 1 g (1) n

 (1) = 0 = 0 is impossible because of gn

0

f (τ ) e−λn (t−τ ) dτ.

(5.60)

206

Inverse Eigenvalue Problems

From partial integration, we observe that t

−λn (t−τ )

f (τ ) e 0

t

1 1 dτ = f (t) − λn λn

f  (τ ) e−λn (t−τ ) dτ ,

0

and this decays as 1/λn . Using this and the asymptotic behavior (5.59a) and (5.59b), we conclude that the series (5.60) converges uniformly in ΩT . For x = 1, the representation (5.60) reduces to U (1, t)

=

∞ 

2

t

gn (1)

n=1

0

t =

f (τ ) e−λn (t−τ ) dτ

f (τ )

∞ 

gn (1)2 e−λn (t−τ ) dτ ,

n=1



0

 =: A(t − τ )

t ∈ [0, T ] .



Changing the orders of integration and summation is justified by Lebesgue’s theorem of dominated convergence. This is seen from the estimate ∞ 

2 −λn s

gn (1) e

≤ c

n=1

∞ 

−n2 π 2 s

e

∞ ≤ c

n=1

e−σ

0

2

π2 s

dσ =

c √ 2 πs

√ and the fact that the function s → 1/ s is integrable in (0, T ]. Such a representation holds for U1 (1, ·) and U2 (1, ·) corresponding to q1 and q2 , respectively. We denote the dependence on q1 and q2 by superscripts (1) and (2), respectively. From U1 (1, ·) = U2 (1, ·), we conclude that t 0 =



(1)

f (τ ) A

(2)

(t − τ ) − A

 (t − τ ) dτ =

0

t

  f (t − τ ) A(1) (τ ) − A(2) (τ ) dτ ;

0

that is, the function w := A(1) − A(2) solves the homogeneous Volterra integral equation of the first kind with kernel f (t − τ ). We differentiate this equation twice and use f (0) = 0 and f  (0) = 0. This yields a Volterra equation of the second kind for w: 

t

f (0) w(t) +

f  (t − s) w(s) ds = 0,

t ∈ [0, T ] .

0

Because Volterra equations of the second kind are uniquely solvable (see Example A.32 of Appendix A.3), this yields w(t) = 0 for all t, that is ∞ ∞    (1) 2 −λ(1) t  (2) 2 −λ(2) t gn (1) e n = gn (1) e n n=1

n=1

for all t ∈ (0, T ) .

5.6

A Parameter Identification Problem

207

(j)

We note that gn (1) > 0 for j = 1, 2 by our normalization. Now we can apply a result from the theory of Dirichlet series (see Lemma 5.26) and conclude that (1) (2) (1) (2) λn = λn and gn (1) = gn (1) for all n ∈ N. Applying the uniqueness result analogous to Theorem 5.23, part (b), for the boundary conditions u(0) = 0 and  u (1) = 0 (see Problem 5.5), we conclude that q1 = q2 . It remains to prove the following lemma: Lemma 5.26 Let λn and μn be strictly increasing sequences that tend to infinity. Let the series ∞ ∞   αn e−λn t and βn e−μn t n=1

n=1

converge for every t ∈ (0, T ] and uniformly on some interval [δ, T ]. Let the limits coincide, that is ∞  n=1

αn e−λn t =

∞ 

βn e−μn t

for all t ∈ (0, T ] .

n=1

If we also assume that αn = 0 and βn = 0 for all n ∈ N, then αn = βn and λn = μn for all n ∈ N. Proof: Assume that λ1 = μ1 or α1 = β1 . Without loss of generality, we can assume that μ1 ≥ λ1 (otherwise, interchange the roles of λn and μn ). Define Cn (t) := αn e−(λn −λ1 )t − βn e−(μn −λ1 )t for t ≥ δ . ∞ By analytic continuation, we conclude that n=1 Cn (t) = 0 for all t ≥ δ and that the series converges uniformly on [δ, ∞). Because C1 (t) = α1 − β1 e−(μ1 −λ1 )t and α1 = β1 or μ1 > λ1 there exist  > 0 and t1 > δ such that |C1 (t)| ≥  for all t ≥ t1 . Choose n0 ∈ N with     n0  <  for all t ≥ t1 .  C (t) n   2 n=1 Then we conclude that       n0    n0    n0        Cn (t) = C1 (t) − Cn (t) ≥ |C1 (t)| −  Cn (t) ≥  2 n=2 n=1 n=1 for all t ≥ t1 . Now we let t tend to infinity. The first finite sum converges to zero, which is a contradiction. Therefore, we have shown that λ1 = μ1 and α1 = β1 . Now we repeat the argument for n = 2, etc. This proves the lemma. 

208

5.7

Inverse Eigenvalue Problems

Numerical Reconstruction Techniques

In this section, we discuss some numerical algorithms for solving the inverse spectral problem which was suggested and tested by W. Rundell, P. Sacks, and others. We follow closely the papers [186, 233, 234]. From now on, we assume knowledge of eigenvalues λn and μn , n ∈ N, of the Sturm–Liouville eigenvalue problems (5.45) or (5.46). It is our aim to determine the unknown function q. Usually, only a finite number of eigenvalues is known. Then one cannot expect to recover the total function q but only “some portion” of it (see (5.62)). The first algorithm we discuss uses the concept of the characteristic function again. For simplicity, we describe the method only for the case where q is known to be an even function; that is, q(1 − x) = q(x). Then we know that only one spectrum suffices to recover q (see Theorem 5.23). Recalling the characteristic function f (λ) = u2 (1, λ, q) for the problem (5.45), the inverse problem can be written as the problem of solving the equations (5.61) u2 (1, λn , q) = 0 for all n ∈ N for the function q. If we know only a finite number, say λn for n = 1, . . . , N , then we assume that q is of the form q(x; a) =

N 

an qn (x) ,

x ∈ [0, 1] ,

(5.62)

n=1

for coefficients a = (a1 , . . . , aN ) ∈ RN and some given linear independent even functions qn . If q is expected to be smooth and periodic, a good choice for qn is qn (x) = cos(2π(n − 1)x), n = 1, . . . , N . Equation (5.61) then reduces to the finite nonlinear system F (a) = 0, where F : RN → RN is defined by Fn (a) := u2 (1, λn , q(·; a))

for a ∈ RN and n = 1, . . . , N .

Therefore, all of the well-developed methods for solving systems of nonlinear equations can be used. For example, Newton’s method  −1  (k) , F a a(k+1) = a(k) − F  a(k)

k = 0, 1, . . . ,

is known to be quadratically convergent if F  (a)−1 is regular. As we know from Section 5.2, Theorem 5.6, the mapping F is continuously Fr´echet differentiable for every a ∈ RN . The computation of the derivative is rather expensive, and in general, it is not known if F  (a) is regular. In [186], it was proven that F  (a) is regular for sufficiently small a and is of triangular form for a = 0. This observation leads to the simplified Newton method of the form  a(k+1) = a(k) − F  (0)−1 F a(k) , k = 0, 1, . . . . For further aspects of this method, we refer to [186].

5.7

Numerical Reconstruction Techniques

209

Before we describe a second algorithm, we observe that from  1 the asymptotic form (5.28) of the eigenvalues, we have an estimate of qˆ = 0 q(x) dx. Writing the differential equation in the form   −un (x) + q(x) − qˆ un (x) = λn − qˆ un (x), 0 ≤ x ≤ 1 , 1 we observe that we can assume without loss of generality that 0 q(x) dx = 0. Now we describe an algorithm that follows the idea of the uniqueness Theorem 5.22. We allow q ∈ L2 (0, 1) to be arbitrary. The algorithm consists of two steps. First, we recover the Cauchy data f = K(1, ·) and g = Kx (1, ·) from the two sets of eigenvalues. Then we suggest Newton-type methods to compute q from these Cauchy data. The starting point is Theorem 5.19 for the case p = 0. We have already formulated this special case in Example 5.20. Therefore, let (λn , un ) be the eigenvalues and eigenfunctions of the eigenvalue problem (5.45) normalized such that un (0) = 1. The eigenvalues λn are assumed to be known. From Example 5.20, we have the representation √ √ x sin λn x sin λn t √ un (x) = + K(x, t) √ dt , λn λn

0 ≤ x ≤ 1,

(5.63)

0

1 where K satisfies (5.53a)–(5.53c) with K(1, 1) = 12 0 q(t) dt = 0. From (5.63) for x = 1, we √ can compute K(1, t) because, by Theorem 5.21, the functions vn (t) = sin λn t form a complete system in L2 (0, 1). When we know only a finite number λ1 , . . . , λN of eigenvalues, we suggest representing K(1, ·) as a finite sum of the form K(1, t) =

N 

ak sin(kπt) ,

k=1

arriving at the finite linear system N  k=1

1 ak

sin(kπt) sin

  λn t dt = − sin λn

for n = 1, . . . , N .

(5.64)

0

The same arguments yield a set of equations for the second boundary condition u (1) + H u(1) = 0 in the form √

√ √ μn cos μn + H sin μn +

1



√ Kx (1, t) + H K(1, t) sin μn t dt = 0 ,

0

where now μn are the corresponding known eigenvalues. The representation Kx (1, t) + H K(1, t) =

N  k=1

bk sin(kπt)

210

Inverse Eigenvalue Problems

leads to the system N  k=1

1 bk

√ √ √ √ sin(kπt) sin μn t dt = − μn cos μn − H sin μn

(5.65)

0

for n = 1, . . . , N . Equations (5.64) and (5.65) are of the same form and we restrict ourselves to the discussion of (5.64). Asymptotically, the matrix A ∈ 1 √ RN ×N defined by Akn = 0 sin(kπt) sin λn t dt is just 12 I. More precisely, from Parseval’s identity (see (A.8) from Theorem A.15 of AppendixA.2) 2 1 ∞ 1     ψ(t) sin(kπt) dt = 1 |ψ(t)|2 dt   2

k=1 0

0

√

we conclude that (set ψ(t) = sin λn t − sin(nπt) for some n ∈ N)  ∞ 1      2  sin(kπt) sin λn t − sin(nπt) dt   k=1 0

=

1 2

1 0

     sin λn t − sin(nπt)2 dt ≤ 1  λn − nπ 2 2

where we used the mean value theorem. The estimate (5.30) yields |λn −n2 π 2 | ≤ c˜q∞ and thus    λn − nπ  ≤ c q∞ , n where c is independent of q and n. From this, we conclude that  ∞ 1 2      2  sin(kπt) sin λn t − sin(nπt) dt ≤ c q2∞ .   2 n2

k=1 0

The matrix A is thus diagonally dominant, and therefore, invertible for sufficiently small q∞ . Numerical experiments have shown that also for “large” values of q the numerical solution of (5.65) does not cause any problems. We are now facing the following inverse problem: Given (approximate values 1 of) the Cauchy data f = K(1, ·) ∈ H00 (0, 1) and g = Kx (1, ·) ∈ L2 (0, 1), 2 compute q ∈ L (0, 1) such that the solution of the Cauchy problem(5.41a)– x (5.41c) for p = 0 and F = 0 assumes the boundary data K(x, x) = 12 0 q(t) dt for x ∈ [0, 1]. An alternative way of formulating the inverse problem is to start with the Goursat problem (5.53a)–(5.53c): Compute q ∈ L2 (0, 1) such that the solution of the initial value problem (5.53a)–(5.53c) has Cauchy data f (t) = K(1, t) and g(t) = Kx (1, t) for t ∈ [0, 1].

5.7

Numerical Reconstruction Techniques

211

We have studied these coupled systems for K and q in Theorem 5.18. Here we apply it for the case where F = 0. It has been shown that the pair (K, r) solves the system ∂ 2 K(x, t) ∂ 2 K(x, t) − − q(x) K(x, t) = 0 in Δ0 , ∂x2 ∂t2

K(x, x) K(x, 0)

= =

1 2

x

0,

r(t) dt ,

0 ≤ x ≤ 1,

0

0 ≤ x ≤ 1,

and

∂ K(1, t) = g(t) for all t ∈ [0, 1] ∂x if and only if w(ξ, η) = K(ξ + η, ξ − η) and r solve the system of integral equations (5.44a) and (5.44b). For this special choice of F , (5.44b) reduces to K(1, t) = f (t) and

1 r(x) = − 2

1

q(s) K(s, 2x − s) ds + g(2x − 1) + f  (2x − 1) ,

(5.66)

x

where we have extended f and g to odd functions on [−1, 1]. Denote by T (q) the expression on the right-hand side of (5.66). For the evaluation of T (q), one has to solve the Cauchy problem (5.41a)–(5.41c) for p = 0. Note that the solution K; that is, the kernel K(y, 2x − y) of the integral operator T , also depends on q. The operator T is therefore nonlinear! The requirement r = q leads to a fixed point equation q = 2T (q) in L2 (0, 1). It was shown in [233] that there exists at most one fixed point q ∈ L∞ (0, 1) of T . Even more, Rundell and Sachs proved that the projected operator PM T is a contraction on the ball BM := q ∈ L∞ (0, 1) : q∞ ≤ M with respect to some weighted L∞ -norms. Here, PM denotes the projection onto BM defined by  q(x), |q(x)| ≤ M, (PM q)(x) = M sign q(x), |q(x)| > M. Also, they showed the effectiveness of the iteration method q (k+1) = 2T (q (k) ) by several numerical examples. We observe that for q (0) = 0 the first iterate q (1) is simply q (1) (x) = 2 g(2x − 1) + 2 f  (2x − 1), x ∈ [0, 1]. We refer to [233] for more details. As suggested earlier, an alternative numerical procedure based on the kernel 2 1 2 function K is to define the operator S from L (0, 1) into H00 (0, 1) × L (0, 1) by S(q) = K(1, ·), Kx (1, ·) , where K solves the Goursat problem (5.53a)–(5.53c) in the weak sense. This operator is well-defined and bounded by Theorem 5.15, 1 (0, 1) and g ∈ L2 (0, 1) are the given Cauchy values K(1, ·) part (c). If f ∈ H00

212

Inverse Eigenvalue Problems

and Kx (1, ·), respectively, then we have to solve the nonlinear equation S(q) = (f, g). Newton’s method does it by the iteration procedure   q (k+1) = q (k) − S  (q (k) )−1 S(q (k) ) − (f, g) , k = 0, 1, . . . . (5.67) For the implementation, one has to compute the Fr´echet derivative of S. Using the Volterra equation (5.40) derived in the proof of Theorem 5.15, it is not diffi-  cult to prove that S is Fr´echet differentiable and that S  (q)r = W (1, ·), Wx (1, ·) , where W solves the inhomogeneous Goursat problem Wxx (x, t) − Wtt (x, t) − q(x) W (x, t) = K(x, t) r(x) W (x, 0) W (x, x)

=

0,

=

1 2

in Δ0 ,

0 ≤ x ≤ 1,

(5.68a) (5.68b)

x 0 ≤ x ≤ 1.

r(t) dt ,

(5.68c)

0

In part (b) of Theorem 5.18, we showed that S  (q) is an isomorphism. We reformulate this result. Theorem 5.27 Let q ∈ L2 (0, 1) and K be the weak solution of (5.53a)–(5.53c). 1 (0, 1) and g ∈ L2 (0, 1), there exists a unique r ∈ L2 (0, 1) and For every f ∈ H00 a weak solution W of (5.68a)–(5.68c) with W (1, ·) = f and Wx (1, ·) = g in the sense of (5.44a), (5.44b); that is, S  (q) is an isomorphism. Implementing Newton’s method is quite expensive because in every step one has to solve a coupled system of the form (5.68a)–(5.68c). Rundell and Sachs suggested a simplified Newton method of the form    q (k+1) = q (k) − S  (0)−1 S q (k) − (f, g) , k = 0, 1, . . . . Because S(0) = 0, we can invert the linear operator S  (0) analytically. In   particular, we have S (0)r = W (1, ·), Wx (1, ·) , where W now solves Wxx (x, t) − Wtt (x, t) = 0 W (x, 0) = 0 ,

1 and W (x, x) = 2

in Δ0 ,

x 0 ≤ x ≤ 1,

r(t) dt, 0

because also K = 0. The solution W of the Cauchy problem Wxx (x, t) − Wtt (x, t) = 0 W (1, t) = f (t),

and

in Δ0 ,

Wx (1, t) = g(t),

0 ≤ t ≤ 1,

is given by 1 W (x, t) = − 2

t+(1−x) 

g(τ ) dτ + t−(1−x)

1  1  f t + (1 − x) + f t − (1 − x) , 2 2

5.8

Problems

213

where we have extended f and g to odd functions again. The solution r of S  (0)r = (f, g) is, therefore, given by r(x) = 2

d W (x, x) = 2 f  (2x − 1) + 2 g(2x − 1) . dx

In this chapter, we have studied only one particular inverse eigenvalue problem. Similar theoretical results and constructive algorithms can be obtained for other inverse spectral problems; see [4, 16]. For an excellent overview, we refer to the lecture notes by W. Rundell [232].

5.8

Problems

5.1 Let q, f ∈ C[0, 1] and q(x) ≥ 0 for all x ∈ [0, 1]. (a) Show that the following boundary value problem on [0, 1] has at most one solution u ∈ C 2 [0, 1]: −u (x) + q(x) u(x) = f (x) ,

u(0) = u(1) = 0 .

(5.69)

(b) Let v1 and v2 be the solutions of the following initial value problems on [0, 1]: −v1 (x) + q(x) v1 (x) = 0 , −v2 (x) + q(x) v2 (x) = 0 ,

v1 (0) = 0 , v1 (0) = 1 , v2 (1) = 0 , v2 (1) = 1 .

Show that the Wronskian v1 v2 −v2 v1 is constant. Define the following function for some a ∈ R:  a v1 (x) v2 (y), 0 ≤ x ≤ y ≤ 1, G(x, y) = a v2 (x) v1 (y), 0 ≤ y < x ≤ 1. Determine a ∈ R such that 1 G(x, y) f (y) dy,

u(x) :=

x ∈ [0, 1] ,

0

solves (5.69). The function G is called Green’s function of the boundary value problem (5.69). (c) Show that the eigenvalue problem −u (x) + q(x) u(x) = λ u(x),

u(0) = u(1) = 0 ,

is equivalent to the eigenvalue problem for the integral equation 1 u(x) = λ

1 G(x, y) u(y) dy, 0

x ∈ [0, 1] .

214

Inverse Eigenvalue Problems Prove Theorem 5.7, parts (a) and (b) by the general spectral theorem (Theorem A.53 of Appendix A.6). (d) How can one treat the case of part (c) when q changes sign?

5.2 Let H ∈ R. Prove that the transcendental equation z cot z + H = 0 has a countable number of zeros zn and that zn = (n + 1/2) π + From this,

H + O(1/n2 ) . (n + 1/2)π

zn2 = (n + 1/2)2 π 2 + 2H + O(1/n)

follows. Hint: Make the substitution z = x + (n + 1/2)π, set ε = 1/(n + 1/2)π, write z cot z +H = 0 in the form f (x, ε) = 0, and apply the implicit function theorem. 5.3 Prove Lemma 5.13. 5.4 Let q ∈ C[0, 1] be real- or complex-valued and λn , gn be the eigenvalues and L2 -normalized eigenfunctions, respectively, corresponding to q and boundary conditions u(0) = 0 and hu (1)+Hu(1) = 0. Show by modifying the proof of Theorem 5.21 that {gn : n ∈ N} is complete in L2 (0, 1). This gives—even for real q—a proof different from the one obtained by applying the general spectral theory. 5.5 Consider the eigenvalue problem on [0, 1]: −u (x) + q(x) u(x) = λ u(x),

u(0) = u (1) = 0 .

By modifying the proof of Theorem 5.22, prove the following uniqueness result for the inverse problem: Let (λn , un ) and (μn , vn ) be the eigenvalues and eigenfunctions corresponding to p and q, respectively. If λn = μn for all n ∈ N and un (1) vn (1) =  for all n ∈ N ,  un (0) vn (0) then p and q coincide.

Chapter 6

An Inverse Problem in Electrical Impedance Tomography 6.1

Introduction

Electrical impedance tomography (EIT) is a medical imaging technique in which an image of the conductivity (or permittivity) of part of the body is determined from electrical surface measurements. Typically, conducting electrodes are attached to the skin of the subject and small alternating currents are applied to some or all of the electrodes. The resulting electrical potentials are measured, and the process may be repeated for numerous different configurations of applied currents. Applications of EIT as an imaging tool can be found in fields such as medicine (monitoring of the lung function or the detection of skin cancer or breast cancer), geophysics (locating of underground deposits, detection of leaks in underground storage tanks), or nondestructive testing (determination of cracks in materials). To derive the EIT model, we start from the time-harmonic Maxwell system in the form curl H + (iωε − γ) E = 0 ,

curl E − iωμ H = 0

in some domain which we take as a cylinder of the form B × R ⊂ R3 with bounded cross-section B ⊂ R2 . Here, ω, ε, γ, and μ denote the frequency, electric permittivity, conductivity, and magnetic permeability, respectively, which are all assumed to be constant along the axis of the cylinder; that is, depend   on x1 and x2 only. We note that the real parts Re exp(−iωt) E(x) and Re exp(−iωt) H(x) are the physically meaningful electric and magnetic field, respectively. For low frequencies ω (i.e., for small (ωμγ) · L2 where L is a typical

© Springer Nature Switzerland AG 2021 A. Kirsch, An Introduction to the Mathematical Theory of Inverse Problems, Applied Mathematical Sciences 120, https://doi.org/10.1007/978-3-030-63343-1 6

215

216

Electrical Impedance Tomography

length scale of B), one can show (see, e.g., [47]) that the Maxwell system is approximated by curl H − γ E = 0 ,

curl E = 0 .

The second equation yields the existence1 of a scalar potential u such that E = −∇u. this into the first equation and taking the divergence yields  Substituting  div γ∇u = 0 in the cylinder. We restrict ourselves to the two-dimensional case and consider the conductivity equation div (γ∇u) = 0

in B .

(6.1)

There are several possibilities for modeling the attachment of the electrodes on the boundary ∂B of B. The simplest of these is the continuum model in which the potential U = u|∂B and the boundary current distribution f = γ ∇u · ν = γ ∂u/∂ν are both given on the boundary ∂B. Here, ν = ν(x) is the unit normal vector at x ∈ ∂B directed into the exterior of B. First, we observe that,2 by the divergence theorem,      ∂u d = div γ∇u dx = γ f d ; 0 = ∂ν B

∂B

∂B

that is, the boundary current distribution f has zero mean. In practice, f (x) is not known for all x ∈ ∂B. One actually knows the currents sent along wires attached to N discrete electrodes that in turn are attached to the boundary ∂B. Therefore, in the gap model one approximates f by assuming that f is constant at the surface of each electrode and zero in the gaps between the electrodes. An even better choice is the complete model . Suppose that fj is the electric current sent through the wire attached to the jth electrode. At the surface Sj of this electrode, the current density satisfies  ∂u d = fj . γ ∂ν Sj

In the gaps between the electrodes, we have γ

 ∂u = 0 in ∂B \ Sj . ∂ν j

If electrochemical effects at the contact of Sj with ∂B are taken into account, the Dirichlet boundary condition u = Uj on Sj is replaced by u + zj γ

∂u = Uj ∂ν

on Sj ,

where zj denotes the surface impedance of the jth electrode. We refer to [20, 145, 146, 252] for a discussion of these electrode models and the well-posedness of the corresponding boundary value problems (for given γ). 1 If

the domain is simply connected. γ, f , and u are smooth enough.

2 Provided

6.2 The Direct Problem

217

In the inverse problem of EIT the conductivity function γ is unknown and has to be determined from simultaneous measurements of the boundary voltages U and current densities f , respectively. In this introductory chapter on EIT, we restrict ourselves to the continuum model as the simplest electrode model. We start with the precise mathematical formulation of the direct and the inverse problem and prove well-posedness of the direct problem: existence, uniqueness, and continuous dependence on both the boundary data f and the conductivity γ. Then we consider the inverse problem of EIT. The question of uniqueness is addressed, and we prove uniqueness of the inverse linearized problem. This problem is interesting also from an historical point of view because the proof, given in Calder´ on’s fundamental paper [38], has influenced research on inverse medium problems monumentally. In the last section, we introduce a technique to determine the support of the contrast γ − γ1 where γ1 denotes the known background conductivity. This factorization method has been developed fairly recently—after publication of the first edition of this monograph—and is a prominent member of a whole class of newly developed methods subsumed under the name Sampling Methods.

6.2

The Direct Problem and the Neumann–Dirichlet Operator

Let B ⊂ R2 be a given bounded domain with boundary ∂B and γ : B → R and f : ∂B → R be given real-valued functions. The direct problem is to determine u such that div (γ∇u) = 0

in B ,

γ

∂u = f ∂ν

on ∂B .

(6.2)

Throughout this chapter, ν = ν(x) again denotes the exterior unit normal  vector at x ∈ ∂B. As mentioned in the introduction, we have to assume that ∂B f d = 0. Therefore, throughout this chapter, we make the following assumptions on B, γ, and f : Assumption 6.1 (a) B ⊂ R2 is a bounded Lipschitz domain3 such that the exterior of B is connected. (b) γ ∈ L∞ (B), and there exists γ0 > 0 such that γ(x) ≥ γ0 for almost all x ∈ B. (c) f ∈ L2 (∂B) where L2 (∂B)

=

2



f ∈ L (∂B) : ∂B

3 For

a definition see, e.g., [191, 161].

f d = 0 .

218

Electrical Impedance Tomography

In this chapter and the following one on scattering theory, we have to use Sobolev spaces as the appropriate functions spaces for the solution. For a general theory on Sobolev spaces, we refer to the standard literature such as [1, 191]. For obvious reasons, we also refer to the monograph [161], Sections 4.1 and 5.1. Also, in Appendix A.5, we introduce and study the Sobolev space H 1 (B) for the particular case of B being the unit disc. At this place we recall only the very basic definition. For any open and bounded set B ⊂ R2 the Sobolev space H 1 (B) is defined as the completion of C 1 (B) with respect to the norm    uH 1 (B) = |∇u|2 + |u|2 dx . B

An important property of Sobolev spaces for Lipschitz domains is the existence of traces; that is, for every u ∈ H 1 (B) the trace u|∂B on ∂B is well-defined and represents an L2 (∂B)-function4 (Theorem 5.10 in [161]). Also, there exists cT > 0 (independent of u) such that u|∂B L2 (∂B) ≤ cT uH 1 (B) for all u ∈ H 1 (B); that is, the trace operator u → u|∂B is bounded. We note that the solution u of (6.2) is only unique up to an additive constant. Therefore, we normalize the solution u ∈ H 1 (B) such that it has vanishing mean on the boundary; that is, u ∈ H1 (B) where

 u d = 0 . (6.3) H1 (B) = u ∈ H 1 (B) : ∂B

The formulation (6.2) of the boundary value problem has to be understood in the variational (or weak) sense. By multiplying the first equation of (6.2) with some test function ψ and using Green’s first formula we arrive at      0 = ψ div γ∇u dx = − γ ∇ψ · ∇u dx + ψ γ∇u · ν d B



= −

 γ ∇ψ · ∇u dx +

B

B

∂B

ψ f d . ∂B

Therefore, we define the variational solution u ∈ H1 (B) of (6.2) by the solution of   γ ∇ψ · ∇u dx = ψ f d for all ψ ∈ H1 (B) . (6.4) B

∂B

Existence and uniqueness follows from the representation theorem due to Riesz (cf. Theorem A.23 of Appendix A.3). Theorem 6.2 Let Assumption 6.1 be satisfied. For every f ∈ L2 (∂B) there exists a unique variational solution u ∈ H1 (B) of (6.2), that is, a solution 4 The

trace is even more regular and belongs to the fractional Sobolev space H 1/2 (∂B).

6.2 The Direct Problem

219

of the variational equation (6.4). Furthermore, there exists a constant c > 0 (independent of f ) such that uH 1 (B) ≤ c f L2 (∂B) . In other words: the operator f → u from L2 (∂B) to H1 (B) is bounded. Proof: We define a new inner product in the space H1 (B) by  (u, v)∗ = γ ∇u · ∇v dx , u, v ∈ H1 (B) . B

 The corresponding norm u∗ = (u, u)∗ is equivalent to the ordinary norm  · H 1 (B) in H1 (B). This follows from Friedrich’s inequality in the form (see Theorem A.50 for the case of B being the unit disc and [1, 191] for more general Lipschitz domains): There exists cF > 0 such that vL2 (B) ≤ cF ∇vL2 (B)

for all v ∈ H1 (B) .

(6.5)

Indeed, from this the equivalence follows inasmuch as γ0 v2H 1 (B) ≤ γ0 ∇v2L2 (B) ≤ v2∗ ≤ γ∞ v2H 1 (B) 1 + c2F

(6.6)

for all v ∈ H1 (B). For fixed f ∈ L2 (∂B) we can interpret the right-hand side of (6.4) as a linear functional F on the space H1 (B); that is, F (ψ) = (f, ψ)L2 (∂B) for ψ ∈ H1 (B). This functional F is bounded by the inequality of CauchySchwarz and the trace theorem (with constant cT > 0) because   F (ψ) ≤ f L2 (∂B) ψL2 (∂B) ≤ cT f L2 (∂B) ψH 1 (B) ≤ c f L2 (∂B) ψ∗  with c = cT (1 + c2F )/γ0 . In particular, F H1 (B)∗ ≤ c f L2 (∂B) , and we can   apply the representation theorem of Riesz in the Hilbert space H1 (B), (·, ·)∗ : there exists a unique u ∈ H1 (B) with (u, ψ)∗ = F (ψ) for all ψ ∈ H1 (B). This is exactly the variational equation (6.4). Furthermore, u∗ = F H1 (B)∗ and thus by (6.6), u2H 1 (B) ≤

1 + c2F 1 + c2F 1 + c2F u2∗ = F 2H1 (B)∗ ≤ c2 f 2L2 (∂B) ; γ0 γ0 γ0

that is, the operator f → u is bounded from L2 (∂B) into H1 (B).



This theorem implies the existence and boundedness of the Neumann–Dirichlet operator. Definition 6.3 The Neumann–Dirichlet operator Λ : L2 (∂B) → L2 (∂B) is defined by Λf = u|∂B , where u ∈ H1 (B) is the uniquely determined variational solution of (6.2); that is, the solution of (6.4).

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Electrical Impedance Tomography

Remark: This operator is bounded by the boundedness of the solution map f → u from L2 (∂B) to H1 (B) and the boundedness of the trace operator from H1 (B) to L2 (∂B). It is even compact because the trace operator is compact from H1 (B) to L2 (∂B). However, we do not make use of this latter property. We show some properties of the Neumann–Dirichlet operator. Theorem 6.4 Let Assumption 6.1 be satisfied. Then the Neumann–Dirichlet map Λ is self-adjoint and positive; that is, (Λf, g)L2 (∂B) = (f, Λg)L2 (∂B) and (Λf, f )L2 (∂B) > 0 for all f, g ∈ L2 (∂B), f = 0. Proof: This follows simply from the definition of Λ and Green’s first identity. Let u, v ∈ H1 (B) be the solutions of (6.4) corresponding to boundary data f and g, respectively. Then, by (6.4) for the pair g, v and the choice ψ = u (note that u|∂B = Λf ),   u g d = γ ∇u · ∇v dx , (Λf, g)L2 (∂B) = ∂B

B

and this term is symmetric with respect to u and v. For f = g this also yields that Λ is positive.  In the following, we write Λγ to indicate the dependence on γ. The next interesting property is a monotonicity result. Theorem 6.5 Let γ1 , γ2 ∈ L∞ (B) with γ1 ≥ γ2 ≥ γ0 a.e. on B. Then Λγ1 ≤ Λγ2 in the sense that     Λγ1 f, f L2 (∂B) ≤ Λγ2 f, f L2 (∂B) for all f ∈ L2 (∂B) . Proof: For fixed f ∈ L2 (∂B) let uj ∈ H1 (B) be the corresponding solution of (6.2) for γj , j = 1, 2. From (6.4) (for γ = γ2 , u = u2 , and ψ = u1 − u2 ) we conclude that     (Λγ1 − Λγ2 )f, f L2 (∂B) = (u1 − u2 ) f d = γ2 (∇u1 − ∇u2 ) · ∇u2 dx ∂B

which proves the result.



B

  γ2 |∇u1 |2 − |∇u2 |2 − |∇(u1 − u2 )|2 dx

=

1 2

≤

1 2

≤

1 2

=

 1 (Λγ1 − Λγ2 )f, f L2 (∂B) 2 

B



B



B

γ2 |∇u1 |2 dx −

1 2

1 γ1 |∇u1 |2 dx − 2



γ2 |∇u2 |2 dx

B



B

γ2 |∇u2 |2 dx

6.3

6.3

The Inverse Problem

221

The Inverse Problem

As described in the introduction, the problem of electrical impedance tomography is to determine (properties of) the conductivity distribution γ from all —or at least a large number of—pairs (f, u|∂B ). Because u|∂B = Λf we can rephrase this problem as follows: Inverse Problem: Determine the conductivity γ from the given Neumann– Dirichlet operator Λγ : L2 (∂B) → L2 (∂B)! As we have seen already in the previous chapter (and this is typical for studying inverse problems), an intensive investigation of the direct problem has to precede the treatment of the inverse problem. In particular, we study the dependence of the Neumann–Dirichlet map on γ. First, we show with an example that the inverse problem of impedance tomography is ill-posed. Example 6.6 Let B = B(0, 1) be the unit disk, qˆ > 0 constant, and R ∈ (0, 1). We define γR ∈ L∞ (B) by 1, R < |x| < 1 , γR (x) = 1 + qˆ , |x| < R . Because γR is piecewise constant the solution u ∈ H1 (B) is a harmonic function in B(0, 1) \ {x : |x| = R} (that is, Δu = 0 in B(0, 1) \ {x : |x| = R}) and satisfies the jump conditions u− = u+ and (1 + qˆ)∂u/∂r|− = ∂u/∂r|+ for |x| = R where v|± denotes the trace of v from the interior (−) and exterior (+) of {x : |x| = R}, respectively. We refer to Problem 6.1 for a justification of this statement. We solve the boundary value problem (6.2) by expanding the boundary data f ∈ L2 (∂B) and the solution u into Fourier series; that is, for  f (ϕ) = fn einϕ , ϕ ∈ [0, 2π] , n=0

we make an ansatz for the solution of (6.2) in the form ⎧  |n| inϕ  ⎪ (bn + cn ) Rr e , ⎪ ⎨ n=0   u(r, ϕ) =  |n|  −|n| inϕ  ⎪ ⎪ bn Rr e + cn Rr , ⎩

r < R, r > R.

n=0

The ansatz already guarantees that u is continuous on the circle r = R. The unknown coefficients bn , cn are determined from the conditions (1+ˆ q ) ∂u/∂r− =   ∂u/∂r + for r = R and ∂u/∂r = f for r = 1. This yields the set of equations (1 + qˆ) (bn + cn ) = bn − cn and bn

|n| − cn |n| R|n| = fn R|n|

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Electrical Impedance Tomography

for all n = 0 which yields explicit formulas for bn and cn . Substituting this into the form of u and taking r = 1 yields (ΛγR f )(ϕ) = u(1, ϕ) =

 α − R2|n| fn einϕ , 2|n| |n| α + R n=0

ϕ ∈ [0, 2π] ,

(6.7)

with α = 1 + 2/ˆ q . We observe that ΛγR is a diagonal operator from L2 (0, 2π) into itself with eigenvalues that behave asymptotically as 1/|n|. Therefore, the natural setting for ΛγR is to consider it as an operator from the Sobolev space −1/2 1/2 H (0, 2π) of order −1/2 into the Sobolev space H (0, 2π) of order 1/2; see Section A.4 of Appendix A. We prefer the setting in L2 (0, 2π) because the more general setting does not give any more insight with respect to the inverse problem. Let Λ1 be the operator with γ = 1 which is given by (Λ1 f )(ϕ) =

 fn einϕ , |n|

ϕ ∈ [0, 2π] .

n=0

We estimate the difference by (ΛγR −

Λ1 )f 2L2 (0,2π)

=

 2   α − R2|n|  |fn |2   2π − 1  α + R2|n|  n2 n=0



R4|n| |fn |2  2 2|n| n2 n=0 α + R

=

8π

≤

4 R4 f 2L2 (0,2π) ; α2

that is,

2R2 ≤ 2R2 α because α ≥ 1. Therefore, we have convergence of ΛγR to Λ1 in the operator norm as R tends to zero. On the other hand, the difference γR − 1∞ = qˆ is constant and does not converge to zero as R tends to zero. This shows clearly that the inverse problem to determine γR from Λ is ill-posed. One can argue that perhaps the sup-norm for γ is not appropriate to measure the error in γ. Our example, however, shows that even if we replace qˆ by a constant qˆR which depends on R such that limR→0 qˆR = ∞ we still have convergence of ΛγR to Λ1 in the operator norm as R tends to zero. Taking, for example, qˆR = qˆ/R3 , we observe that γR − 1Lp (B) → ∞ as R tends to zero for arbitrary p ≥ 1, and the problem of impedance tomography is also ill-posed with respect to any Lp -norm. ΛγR − Λ1 L(L2 (0,2π)) ≤

A fundamental question for every inverse problem is the question of uniqueness: is the information—at least in principle—sufficient to determine the

6.3

The Inverse Problem

223

unknown quantity? Therefore, in electrical impedance tomography, we ask: does the knowledge of the Neumann–Dirichlet operator Λ determine the conductivity γ uniquely or is it possible that two different γ correspond to the same Λ? In full generality, this fundamental question was not answered until 2006 by K. Astala and L. P¨ aiv¨ arinta in [10]. We state the result without proof. Theorem 6.7 Let γ1 , γ2 ∈ L∞ (B) with γj (x) ≥ γ0 for j = 1, 2 and almost all x ∈ B. We denote the corresponding Neumann–Dirichlet operators by Λ1 and Λ2 , respectively. If Λ1 = Λ2 then γ1 = γ2 in B. Instead of proving this theorem which uses refined arguments from complex analysis, we consider the linearized problem. Therefore, writing Λ(γ) instead of Λγ to indicate the dependence on γ, we consider the linear problem Λ(γ) + Λ (γ)q = Λmeas , (6.8)   where Λ (γ) : L∞ (B) → L L2 (∂B) denotes the Fr´echet  derivative of the non- linear operator γ → Λ(γ) from L∞ (B) to L L2 (∂B) at γ. Here, L L2 (∂B) denotes again the space of all linear and bounded operators from L2 (∂B) into  itself equipped with the operator norm. The right-hand side Λmeas ∈ L L2 (∂B) is given (“measured”), and the contrast q ∈ L∞ (B) has to be determined. Theorem 6.8   Let U ⊂ L∞ (B) be given by U = γ ∈ L∞ (B) : γ ≥ γ0 a.e. on B .   (a) The mapping γ → Λ(γ) from U to L L2 (∂B) is Lipschitz continuous.   (b) The mapping γ → Λ(γ) from U to L L2 (∂B) is Fr´echet differentiable.  The Fr´echet U in the direction q ∈ L∞ (B) is given  derivative Λ (γ) at γ ∈  1 by Λ (γ)q f = v|∂B where v ∈ H (B) solves     div γ∇v = − div q∇u in B ,

γ

∂u ∂v = −q on ∂B , ∂ν ∂ν

(6.9)

and u ∈ H1 (B) solves (6.2) with data γ ∈ U and f ∈ L2 (∂B). The solution of (6.9) is again understood in the weak sense; that is,   γ ∇ψ · ∇v dx = − q ∇ψ · ∇u dx for all ψ ∈ H1 (B) . (6.10) B

B

Proof: (a) Let γ1 , γ2 ∈ U , f ∈ L2 (∂B), and u1 , u2 ∈ H1 (B) be the corresponding weak solutions of (6.2). Taking the difference of (6.4) for the triples (γ1 , u1 , f ) and (γ2 , u2 , f ) yields   γ1 ∇(u1 − u2 ) · ∇ψ dx = (γ2 − γ1 ) ∇u2 · ∇ψ dx for all ψ ∈ H1 (B) . B

B

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Electrical Impedance Tomography

With ψ = u1 − u2 and the lower bound γ0 ≤ γ1 this yields  2  γ0 ∇(u1 − u2 )2L2 (B) ≤ γ1 ∇(u1 − u2 ) dx B



(γ2 − γ1 ) ∇u2 · ∇(u1 − u2 ) dx

= B

≤

γ1 − γ2 ∞ ∇(u1 − u2 )L2 (B) ∇u2 L2 (B) ;

that is, there exists a constant c1 > 0 (independent of γ1 , γ2 ) with ∇(u1 − u2 )L2 (B)

1 γ1 − γ2 ∞ ∇u2 L2 (B) γ0 ≤ c1 γ1 − γ2 ∞ f L2 (∂B) , ≤

(6.11)

where we use the boundedness of the mapping f → u2 (see Theorem 6.2). Now we use the trace theorem and (6.6) to conclude that    Λ(γ1 )f − Λ(γ2 )f L2 (∂B) = (u1 − u2 )∂B L2 (∂B) ≤ c2 γ1 − γ2 ∞ f L2 (∂B) ; that is, Λ(γ1 ) − Λ(γ2 )L(L2 (∂B)) ≤ c2 γ1 − γ2 ∞ which proves part (a). (b) Let γ ∈ U and q ∈ L∞ (B) such that q∞ ≤ γ0 /2. Then γ + q ≥ γ0 /2 a.e. on B. Let u, uq ∈ H1 (B) correspond to γ and γ + q, respectively, and boundary data f . Subtraction of (6.4) for the triple (γ, u, f ) and (6.10) from (6.4) for (γ + q, uq , f ) yields   γ ∇(uq − u − v) · ∇ψ dx = q ∇(u − uq ) · ∇ψ dx for all ψ ∈ H1 (B) . B

B

Taking ψ = uq − u − v yields as in part (a) an estimate of the form ∇(uq − u − v)L2 (B) ≤

1 q∞ ∇(u − uq )L2 (B) . γ0

Now we use (6.11) (with u1 = u and u2 = uq ) to conclude that ∇(uq − u − v)L2 (B) ≤

c1 q2∞ f L2 (∂B) . γ0

Again by the trace theorem and (6.6) this yields       Λ(γ + q)f − Λ(γ)f − Λ (γ)q f L2 (∂B) = (uq − u − v)  2 ∂B L (∂B) ≤ which proves part (b).



c q2∞ f L2 (∂B) ,

6.3

The Inverse Problem

225

We now show that, for any given constant background medium γ, the linearized inverse problem of electrical impedance tomography (6.8) has at most one solution; that is, the Fr´echet derivative is one-to-one. As already mentioned in the introduction, this proof is due to Calder´ on (see [38]) and has “opened the door” to many uniqueness results in tomography and scattering theory. We come back to this method in the next chapter where we prove uniqueness of an inverse scattering problem by this method.  ∞ Theorem 6.9 Let γ be constant. Then the Fr´echet derivative Λ (γ) : L (B) →  2 L L (∂B) is one-to-one.

Proof: First we note that we can assume without loss of generality that γ = 1.  Let q ∈ L∞ (B) such that Λ (γ)q = 0; that is, Λ (γ)q f = 0 for all f ∈ L2 (∂B). The proof consists of two parts. First, we show that q is orthogonal to all products of two gradients of harmonic functions. Then, in the second part, by choosing special harmonic functions we show that the Fourier transform of q vanishes. Let u1 ∈ C 2 (B) be any harmonic function; that is, Δu1 = 0 in B. Define f ∈ L2 (∂B) by f = ∂u1 /∂ν on ∂B. Then u1 is the solution of (6.2) with Neumann boundary data f . We denote by v1 ∈ H1 (B) the corresponding solution of (6.10); that is, 

 ∇ψ · ∇v1 dx = − B

q ∇ψ · ∇u1 dx

for all ψ ∈ H1 (B) .

B

Now we take a second arbitrary harmonic function u2 ∈ C 2 (B) and set ψ = u2 in the previous equation. This yields 

 q ∇u2 · ∇u1 dx = −

B

 ∇u2 · ∇v1 dx = −

B

∂B

v1

∂u2 d ∂ν

  by Green’s first theorem. Now we note that v1 |∂B = Λ (γ)q f = 0. Therefore, we conclude that the right-hand side vanishes; that is  q ∇u2 · ∇u1 dx = 0 for all harmonic functions u1 , u2 ∈ C 2 (B) . (6.12) B

So far, we considered real-valued functions u1 and u2 . By taking the real and imaginary parts, we can also allow complex-valued harmonic functions for u1 and u2 . Now we fix any y ∈ R2 with y = 0. Let y ⊥ ∈ R2 be a vector (unique up z ± ∈ C2 to sign) with y · y ⊥ = 0 and |y| = |y ⊥ |. Define the complex  vectors ± 2 1 ± ⊥ ± ± by z = 2 (iy ± y ). Then one computes that z · z = j=12 (zj ) = 0 and  z + · z − = j=12 zj+ zj− = − 12 |y|2 and z + + z − = iy. From this, we observe that

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Electrical Impedance Tomography

the functions u± (x) = exp(z ± · x), x ∈ R2 , are harmonic in all of R2 . Therefore, substituting u+ and u− into (6.12) yields  0 = B

q ∇u+ ·∇u− dx = z + ·z −

 B

q(x) e(z

+

+z − )·x

1 dx = − |y|2 2



q(x) eiy·x dx.

B

From this, we conclude that the Fourier transform of q (extended by zero in the exterior of B) vanishes on R2 \ {0}, and thus also q itself. This ends the proof. 

6.4

The Factorization Method

In this section, we consider the full nonlinear problem but restrict ourselves to the more modest problem to determine only the shape of the region D, where γ differs from the known background medium which we assume to be homogeneous with conductivity 1. We sharpen the assumption on γ of Assumption 6.1. Assumption 6.10 In addition to Assumption 6.1, let there exist finitely many domains Dj , j = 1, . . . , m, such that Dj ⊂ B and Dj ∩ Dk = ∅ forj = k and m such that the complement B \ D of the closure of the union D = j=1 Dj is connected. Every domain Dj is assumed to satisfy the exterior cone condition (see. e.g., [103]); that is, for every z ∈ ∂Dj there exists a set C (part of a cone) of the form

2 ˆ x > 1 − δ, 0 < |x| < ε0 C = z + x∈R :θ· |x| ˆ = 1, such that C ∩ Dj = ∅. for some ε0 , δ > 0 and θˆ ∈ R2 with |θ| Furthermore, there exists q0 > 0 such that γ = 1 on B \ D and γ ≥ 1 + q0 on D. We define the contrast q ∈ L∞ (B) by q = γ − 1 and note that D is the support of q. It is not difficult to show that every Lipschitz domain D satisfies the exterior cone condition (see Problem 6.6). The inverse problem of this section is to determine the shape of D from the Neumann–Dirichlet operator Λ. In the following, we use the relative data Λ − Λ1 where Λ1 : L2 (∂B) → corresponds to the known background medium; that is, to γ = 1. The information that Λ − Λ1 does not vanish simply means that the background is perturbed by some contrast q = γ − 1. In the factorization method, we develop a criterion to decide whether or not a given point z ∈ B belongs to D. The idea is then to take a fine grid in B and to check this criterion for every grid point z. This provides a pixel-based picture of D. L2 (∂B)

6.4

The Factorization Method

227

We recall that Λf = u|∂B and Λ1 f = u1 |∂B , where u, u1 ∈ H1 (B) solve  (1 + q) ∇u · ∇ψ dx = (f, ψ)L2 (∂B) for all ψ ∈ H1 (B) , (6.13) B

 ∇u1 · ∇ψ dx

=

(f, ψ)L2 (∂B)

for all ψ ∈ H1 (B) .

(6.14)

B

For the difference, we have (Λ1 − Λ)f = (u1 − u)|∂B , and u1 − u ∈ H1 (B) satisfies the variational equation   (1 + q) ∇(u1 − u) · ∇ψ dx = q ∇u1 · ∇ψ dx for all ψ ∈ H1 (B) . (6.15) B

D

It is the aim to factorize the operator Λ1 − Λ in the form Λ1 − Λ = A∗ T A , where the operators A : L2 (∂B) → L2 (D)2 and T : L2 (D)2 → L2 (D)2 are defined as follows:5 • Af = ∇u1 |D , where u1 ∈ H1 (B) solves the variational equation (6.14), and • T h = q(h − ∇w) where w ∈ H1 (B) solves the variational equation   (1 + q) ∇w · ∇ψ dx = q h · ∇ψ dx for all ψ ∈ H1 (B) . (6.16) B

D

We note that the solution w of (6.16) exists and is unique. This is seen by the representation theorem A.23 of Riesz because the right-hand side again defines  a linear and bounded functional F (ψ) = D q h·∇ψ dx on H1 (B). The left-hand side of (6.16) is again the inner product (w, ψ)∗ . The classical interpretation of the variational equation (under the assumption that all functions are sufficiently smooth) can again be seen from Green’s first theorem, applied in D and in B\D. Indeed, in this case (6.16) is equivalent to       ψ div (1 + q) ∇w − q h dx − ψ ν · (1 + q) ∇w − q h d 0 = D



 ψ Δw dx −

+ B\D

∂(B\D)

∂D

∂w d ψ ∂ν

for all ψ. This yields   div (1 + q) ∇w − q h = 0 in D ,

Δw = 0 in B \ D ,

5 Here, L2 (D)2 denotes the space of vector-valued functions D → R2 such that both components are in L2 (D).

228

Electrical Impedance Tomography   ∂w   ν · (1 + q) ∇w − q h − = on ∂D , ∂ν +

and



∂w = 0 on ∂B ; ∂ν

that is,    ∂w   ν · (1 + q) ∇w − − = q|− ν · h on ∂D , ∂ν +

∂w = 0 on ∂B . ∂ν

Theorem 6.11 Let the operators A : L2 (∂B) → L2 (D)2 and T : L2 (D)2 → L2 (D)2 be defined as above by (6.14) and (6.16), respectively. Then Λ1 − Λ = A∗ T A .

(6.17)

Proof: We define the auxiliary operator H : L2 (D)2 → L2 (∂B) by Hh = w|∂B where w ∈ H1 (B) solves (6.16). Obviously, we conclude from (6.15) that Λ1 − Λ = HA. We determine the adjoint A∗ : L2 (D)2 → L2 (∂B) of A and prove that ∗ A h = v|∂B where v ∈ H1 (B) solves the variational equation   ∇v · ∇ψ dx = h · ∇ψ dx for all ψ ∈ H1 (B) (6.18) B

D

and even for all ψ ∈ H 1 (B) because it obviously holds for constants. The solution v exists and is unique by the same arguments as above. Again, by applying Green’s theorem we note that v is the variational solution of the boundary value problem ∂v div h in D , = 0 on ∂B , (6.19a) Δv = 0 in B \ D , ∂ν ∂ ∂ v|− − v|+ = ν · h on ∂D . v|− = v|+ on ∂D , (6.19b) ∂ν ∂ν To prove the representation of A∗ h, we conclude from the definition of A, equation (6.18) for ψ = u1 , and (6.14) that   ∇u1 · h dx = ∇u1 · ∇v dx = (f, v)L2 (∂B) , (Af, h)L2 (D)2 = D

B

and thus v|∂B is indeed the value of the adjoint A∗ h. Now it remains to show that H = A∗ T . Let h ∈ L2 (D)2 and w ∈ H1 (B) solve (6.16). Then Hh = w|∂B . We rewrite (6.16) as   ∇w · ∇ψ dx = q (h − ∇w) · ∇ψ dx for all ψ ∈ H1 (B) . (6.20) B

D

  The comparison with (6.18) yields A∗ q(h − ∇w) = w|∂B = Hh; that is,  A∗ T = H. Substituting this into Λ1 − Λ = HA yields the assertion. Properties of the operators A and T are listed in the following theorem:

6.4

The Factorization Method

229

Theorem 6.12 The operator A : L2 (∂B) → L2 (D)2 is compact, and the operator T : L2 (D)2 → L2 (D)2 is self-adjoint and coercive (T h, h)L2 (D)2 ≥ c h2L2 (D)2   where c = q0 1 − q0 /(1 + q0 ) > 0.

for all h ∈ L2 (D)2 ,

(6.21)

Proof: (i) For smooth functions u1 ∈ C 2 (B) with Δu1 = 0 in B and ∂u1 /∂ν = f on ∂B the following representation formula holds (see [53] or Theorem 7.16 for the case of the three-dimensional Helmholtz equation).    ∂ ∂u1 (y) − u1 (y) Φ(x, y) d(y) u1 (x) = Φ(x, y) ∂ν ∂ν(y) ∂B    ∂ = Φ(x, y) d(y) , x ∈ B , Φ(x, y) f (y) − (Λ1 f )(y) ∂ν(y) ∂B

where Φ denotes the fundamental solution of the Laplace equation in R2 ; that is, 1 ln |x − y| , x = y . Φ(x, y) = − 2π  We can write ∇u1 in D in the form ∇u1 D = K1 f − K2 Λ1 f where the operators K1 , K2 : L2 (∂B) → L2 (D)2 , defined by  (K1 f )(x) = ∇ Φ(x, y) f (y) d(y) , x ∈ D , ∂B



(K2 g)(x)

= ∇ ∂B

g(y)

∂ Φ(x, y) d(y) , ∂ν(y)

x ∈ D,

are compact as integral operators on bounded regions of integration with smooth kernels. (Note that D ⊂ B.) The representation A = K1 − K2 Λ1 holds by a density argument (see Theorem A.30). Therefore, also A is compact. (ii) Let h1 , h2 ∈ L2 (D)2 with corresponding solutions w1 , w2 ∈ H1 (B) of (6.16). Then, with (6.16) for h2 , w2 and ψ = w1 :  q (h1 − ∇w1 ) · h2 dx (T h1 , h2 )L2 (D)2 = D





q h1 · h2 dx −

= D

q ∇w1 · h2 dx D





q h1 · h2 dx −

= D

(1 + q) ∇w1 · ∇w2 dx . B

This expression is symmetric with respect to h1 and h2 . Therefore, T is selfadjoint.

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Electrical Impedance Tomography

For h ∈ L2 (D)2 and corresponding solution w ∈ H1 (B) of (6.16), we conclude that   (T h, h)L2 (D)2 = q |h − ∇w|2 dx + q (h − ∇w) · ∇w dx D

D





2

q |h − ∇w| dx +

= D



≥

|∇w|2 dx

(with the help of (6.20))

B

  q0 |h|2 − 2 q0 h · ∇w + (1 + q0 ) |∇w|2 dx

D

     1 + q0 ∇w − √ q0 =  1 + q0 D  q0 ≥ q0 1 − h2L2 (D)2 . 1 + q0

2   q0 h + q0 1 − 1 + q0

! |h|

2

dx

 From this result and the factorization (6.17), we note that Λ1 −Λ is compact, self-adjoint (this follows already from Theorem 6.4), and nonnegative. Now we derive the binary criterion on a point z ∈ B to decide whether or not this point belongs to D. First, for every point z ∈ B we define a particular function G(·, z) such that ΔG(·, z) = 0 in B \ {z} and ∂G(·, z)/∂ν = 0 on ∂B such that G(x, z) becomes singular as x tends to z. We construct G from the Green’s function N for Δ in B with respect to the Neumann boundary conditions. ˜ (x, z) where We make an ansatz for N in the form N (x, z) = Φ(x, z) − N again 1 ln |x − z| , x = z , Φ(x, z) = − 2π ˜ (·, z) ∈ is the fundamental solution of the Laplace equation in R2 and determine N 1 H (B) as the unique solution of the Neumann problem ˜ ∂N ∂Φ 1 (·, z) = (·, z) + on ∂B . ∂ν ∂ν |∂B|    We note that the solution exists because ∂Φ(·, z)/∂ν + 1/|∂B| d = 0. This ˜ (·, z) = 0 in B ΔN

and

∂B

is seen by Green’s first theorem in the region B \ B(z, ε):   ∂Φ ∂Φ (·, z) d = (x, z) d(x) ∂ν ∂ν ∂B

|x−z|=ε

= −

1 2π



|x−z|=ε

x−z x−z d(x) = −1 . · 2 |x − z| |x − z|

6.4

The Factorization Method

231

˜ is the Green’s function in B with respect to the Neumann Then N = Φ − N boundary conditions; that is, N satisfies ΔN (·, z) = 0 in B \ {z} and

1 ∂N (·, z) = − on ∂B . ∂ν |∂B|

˜ (·, z) on the parameter From the differentiable dependence of the solution N 2 z ∈ B, we conclude that, for any fixed a ∈ R with |a| = 1, the function G(·, z) = a · ∇z N (·, z) satisfies ΔG(·, z) = 0 in B \ {z} and

∂G (·, z) = 0 on ∂B . ∂ν

(6.22)

The function G(·, z) has the following desired properties. Lemma 6.13 Let z ∈ B, ε0 > 0, θ ∈ [0, 2π], and δ > 0 be kept fixed. For ε ∈ [0, ε0) define the set (part of a cone, compare with Assumption 6.10 for θ θˆ = cos sin θ ) 

cos t Cε = z + r : ε < r < ε0 , |θ − t| < arccos(1 − δ) sin t with vertex in z. Let ε0 be so small such that C0 ⊂ B. Then lim G(·, z)L2 (Cε ) = ∞ .

ε→0

˜ (·, z) it is sufficient to consider only the part a · Proof: By the smoothness of N  t ∇z ln |x−z|. Using polar coordinates for x with respect to z (i.e., x = z+r cos sin t ), s , we have with η = arccos(1 − δ) and the representation of a as a = cos sin s 

  a · ∇z ln |x − z|2 dx

  =

Cε

2 ε0 θ+η  2 (z − x) · a r cos2 (s − t) dx = r dt dr |x − z|4 r4 ε θ−η

Cε θ+η 

cos2 (s − t) dt

= θ−η

"

Therefore, G(·, z)L2 (Cε ) ≥

& c ln

#$

= c

%

ε0 ε

ε0 1 dr = c ln . r ε

ε0 ˜ (·, z)L2 (C ) → ∞ − a · ∇z N 0 ε

for ε → 0 .

 We observe that the functions φz (x) = G(·, z)|∂B are traces of harmonic functions in B \ {z} with vanishing normal derivatives on ∂B. Comparing this with the classical formulation (6.19a) (6.19b) of the adjoint A∗ of A it seems to be plausible that the “source region” D of (6.19a), (6.19b) can be determined by moving the source point z in φz . This is confirmed in the following theorem.

232

Electrical Impedance Tomography

Theorem 6.14 Let Assumptions 6.10 hold and let a ∈ R2 with |a| = 1 be fixed. For every z ∈ B define φz ∈ L2 (∂B) by φz (x) = G(x, z) = a · ∇z N (x, z) ,

x ∈ ∂B ,

(6.23)

where N denotes the Green’s function with respect to the Neumann boundary condition. Then (6.24) z ∈ D ⇐⇒ φz ∈ R(A∗ ) , where A∗ : L2 (D)2 → L2 (∂B) is the adjoint of A, given by (6.18), and R(A∗ ) its range. Proof: First let z ∈ D. Choose a disc B[z, ε] = {x ∈ R2 : |x − z| ≤ ε} with center z and radius ε > 0 such that B[z, ε] ⊂ D. Furthermore, choose a function ϕ ∈ C ∞ (R2 ) such that ϕ(x) = 0 for |x − z| ≤ ε/2 and ϕ(x) = 1 for |x − z| ≥ ε and set w(x) = ϕ(x)G(x, z) for x ∈ B. Then w ∈ H1 (B) and w = G(·, z) in B \ D, thus w|∂B = φz . Next, we determine u ∈ H1 (D) as a solution of Δu = Δw in D, ∂u/∂ν = 0 on ∂D; that is, in weak form    ∂ G(·, z) d , ψ ∈ H1 (D) , ∇u · ∇ψ dx = ∇w · ∇ψ dx − ψ ∂ν D

D

∂D

because ∂w/∂ν = ∂G(·, z)/∂ν on ∂D. Again, the solution exists and is unique. Application of Green’s first theorem in B \ D yields   ∂ ∂ G(·, z) d = G(·, z) d = 0 . ∂ν ∂ν ∂D

∂B

Therefore, the previous variational equation holds also for constants and thus for all ψ ∈ H 1 (D). Now let ψ ∈ H1 (B) be a test function on B. Then    ∂ G(·, z) d ∇u · ∇ψ dx = ∇w · ∇ψ dx − ψ ∂ν D D ∂D    = ∇w · ∇ψ dx + ∇G(·, z) · ∇ψ dx = ∇w · ∇ψ dx . D

B

B\D

Therefore, the definition h = ∇u in D yields A∗ h = w|∂B = φz and thus φz ∈ R(A∗ ). Now we prove the opposite direction. Let z ∈ / D. We have to show that φz is not contained in the range of A∗ and assume, on the contrary, that φz = A∗ h for some h ∈ L2 (D)2 . Let v ∈ H1 (B) be the corresponding solution of (6.18). Therefore, the function w = v − G(·, z) vanishes on ∂B and solves the following equations in the weak form Δw = 0 in B \ Dε (z) ,

∂w = 0 on ∂B , ∂ν

6.4

The Factorization Method

233

for every ε > 0 such that Dε (z) := D ∪ B(z, ε) ⊂ B; that is,  ∇w · ∇ψ dx = 0 B\Dε (z)

  for all ψ ∈ H 1 B \ Dε (z) withψ = 0 on ∂D  ε (z). We extend w by zero into the exterior of B. Then w ∈ H 1 R2 \ Dε (z) because w = 0 on ∂B 6 and  ∇w · ∇ψ dx = 0 R2 \Dε (z)

  for all ψ ∈ H 1 R2 \Dε (z) which vanish on ∂Dε (z). This is the variational form small ε > 0 we of Δw = 0 in R2 \ Dε (z). Since this holds for all sufficiently  conclude that Δw = 0 in the exterior Ω := R2 \ D ∪ {z} of D ∪ {z}. Now we use without proof7 that w is analytic in this set Ω and thus satisfies the unique continuation principle, see, e.g., Theorem 4.39 of [161]. Therefore, because it vanishes in the exterior of B it has to vanish in all of the connected set Ω. (Here we make  assumption that B \ D is connected.) Therefore, v = G(·, z)  use of the in B \ D ∪ {z} . The point z can either be on the boundary ∂D or of D.   in the  exterior t : 0 < r < In either case there is a cone C0 of the form C0 = z + r cos sin t  ε0 , |θ − t| < δ with C0 ⊂ B \ D. (Here we use the fact that every component of D satisfies the exterior cone condition.) It is v|C0 ∈ L2 (C0 ) because even v ∈ H1 (B). However, Lemma 6.13 yields that G(·, z)L2 (Cε ) → ∞ for ε → 0   t  where Cε = z + r cos sin t : ε < r < ε0 , |θ − t| < δ . This is a contradiction  because v = G(·, z) in C0 and ends the proof. Therefore, we have shown an explicit characterization of the unknown domain D by the range of the operator A∗ . This operator, however, is also unknown: only Λ1 − Λ is known! The operators A∗ and Λ1 − Λ are connected by the factorization Λ1 − Λ = A∗ T A. We can easily derive a second factorization of Λ1 − Λ. The operator Λ1 − Λ is self-adjoint and compact as an operator from L2 (∂B) into itself. Therefore, there exists a spectral decomposition of the form (Λ1 − Λ)f =

∞ 

λj (f, ψj )L2 (∂B) ψj ,

n=1

where λj ∈ R denote the eigenvalues and ψj ∈ L2 (∂B) the corresponding orthonormal eigenfunctions of Λ1 − Λ (see Theorem A.53 of Appendix A.6). Furthermore, from the factorization and the coercivity of T it follows that   is not quite obvious that the extension is in H 1 R2 \ Dε (z) , see, e.g., Corollary 5.13 in [161] 7 see again [161], Theorem 4.38 and Corollary 3.4 6 It

234

Electrical Impedance Tomography

  (Λ1 − Λ)f, f L2 (∂B) ≥ 0 for all f ∈ L2 (∂B). This implies λj ≥ 0 for all j. Therefore, we can define Wf =

∞   λj (f, ψj )L2 (∂B) ψj , n=1

and have a second factorization in the form W W = Λ1 −Λ. We write (Λ1 −Λ)1/2 for W . The operator (Λ1 − Λ)1/2 is also self-adjoint, and we have (Λ1 − Λ)1/2 (Λ1 − Λ)1/2 = Λ1 − Λ = A∗ T A .

(6.25)

We show that the ranges of (Λ1 − Λ)1/2 and A∗ coincide.8 This follows directly from the following functional analytic result: Lemma 6.15 Let X and Y be Hilbert spaces, B : X → X, A : X → Y , and T : Y → Y linear and bounded such that B = A∗ T A. Furthermore, let T be self-adjoint and coercive; that is, there exists c > 0 such that (T y, y)Y ≥ cy2Y for all y ∈ Y . Then, for any φ ∈ X, φ = 0,   φ ∈ R(A∗ ) ⇐⇒ inf (Bx, x)X : x ∈ X , (x, φ)X = 1 > 0 . Proof: (i) First, let φ = A∗ y ∈ R(A∗ ) for some y ∈ Y . Then y = 0, and we estimate for arbitrary x ∈ X with (x, φ)X = 1: (Bx, x)X

= = =

(A∗ T Ax, x)X = (T Ax, Ax)Y ≥ cAx2Y 2 c c  (Ax, y)Y  Ax2Y y2Y ≥ 2 2 yY yY  2 c  c  c 2 (x, A∗ y)X  = (x, φ)X  = . 2 2 yY yY y2Y

Therefore, we have found a positive lower bound for the infimum. (ii) Second, let φ ∈ / R(A∗ ). Define the closed subspace   V = x ∈ X : (φ, x)X = 0 = {φ}⊥ . We show that the image A(V ) is dense in the closure of the range of A. Indeed, let y ∈ closure(R(A)) such that y ⊥ Ax for all x ∈ V ; that is, 0 = (Ax, y)Y = / R(A∗ ) we (x, A∗ y) for all x ∈ V ; that is, A∗ y ∈ V ⊥ = span{φ}. Because φ ∈ ∗ ∗ conclude that A y = 0. Therefore, y ∈ closure(R(A)) ∩ N (A ). This yields y = 0.9 Therefore, A(V ) is dense in closure(R(A)). Because Aφ/φ2X is in the xn → −Aφ/φ2X . We range of A there exists a sequence x ˜n ∈ V such that A˜ 2 ˜n + φ/φX . Then (xn , φ)X = 1 and Axn → 0 for n → ∞, and define xn := x we estimate (Bxn , xn )X = (T Axn , Axn )Y ≤ T L(Y ) Axn 2Y −→ 0 , 8 This

n → ∞,

is also known as Douglas’ Lemma, see [75]. a sequence (xj ) in X such that Axj → y. Then 0 = (A∗ y, xj )X = (y, Axj )Y → (y, y)Y ; that is, y = 0. 9 Take

6.4

The Factorization Method

235

  and thus inf (Bx, x)X : x ∈ X , (x, φ)X = 1 = 0.



We apply this result to both of the factorizations of (6.25). In both cases, B = Λ1 − Λ and X = L2 (∂B). First, we set Y = L2 (D)2 and A : L2 (∂B) → L2 (D)2 and T : L2 (D)2 → L2 (D)2 as in the second factorization of (6.25). Because T is self-adjoint and coercive we conclude for any φ ∈ L2 (∂B), φ = 0, that    φ ∈ R(A∗ ) ⇔ inf (Λ1 −Λ)f, f L2 (∂B) : f ∈ L2 (∂B) , (f, φ)L2 (∂B) = 1 > 0 . Second, we consider the first factorization of (6.25) with T being the identity. For φ ∈ L2 (∂B), φ = 0, we conclude that      φ ∈ R (Λ1 − Λ)1/2 ⇔ inf (Λ1 − Λ)f, f L2 (∂B) : (f, φ)L2 (∂B) = 1 > 0 . The right-hand sides of the characterizations only depend on Λ1 − Λ, therefore, we conclude that   (6.26) R (Λ1 − Λ)1/2 = R(A∗ ) . Application of Theorem 6.14 yields the main result of the factorization method: Theorem 6.16 Let Assumptions 6.10 be satisfied. For fixed a ∈ R2 with a = 0 and every z ∈ B let φz ∈ L2 (∂B) be defined by (6.23); that is, φz (x) = a · ∇z N (x, z), x ∈ ∂B, where N denotes the Green’s function for Δ with respect to the Neumann boundary conditions. Then   (6.27) z ∈ D ⇐⇒ φz ∈ R (Λ1 − Λ)1/2 . We now rewrite the right-hand side with Picard’s Theorem A.58 of Appendix A.6. First, we show injectivity of the operator Λ1 − Λ. Theorem 6.17 The operator Λ1 − Λ is one-to-one. Proof: From   (Λ1 − Λ)f, f L2 (∂B)

=

(A∗ T Af, f )L2 (∂B) = (T Af, Af )L2 (D)2

≥

cAf 2L2 (D)2

for f ∈ L2 (∂B)

it suffices to prove injectivity of A. Let Af = ∇u1 |D = 0 where u1 ∈ H1 (B) denotes the weak solution of Δu1 = 0 in B and ∂u1 /∂ν = f on ∂B. Therefore, ∇u1 is constant in every component of D. Without proof, we use again the regularity result that u1 is analytic in B. The derivatives vj = ∂u1 /∂xj are solutions of Δvj = 0 in B and vj = 0 in D. The unique continuation property  yields vj = ∂u1 /∂xj = 0 in all of B and thus f = 0. Therefore, the operator Λ1 − Λ is self-adjoint, compact, one-to-one, and all eigenvalues are positive. Let {λj , ψj } be an eigensystem of Λ1 − Λ; that

236

Electrical Impedance Tomography

is, λj > 0 are the eigenvalues of Λ1 − Λ and ψj ∈ L2 (∂B) are the corresponding orthonormal eigenfunctions (see Theorem A.53 of Appendix A.6). The  set {ψj : j = 1, 2, . . .} is complete by the spectral theorem. Therefore, { λj , ψj , ψj } is a singular system of the operator (Λ1 − Λ)1/2 . Application of Picard’s Theorem A.58 of Appendix A.6 yields Theorem 6.18 Let Assumptions 6.10 be satisfied. For fixed a ∈ R2 with a = 0 and for every z ∈ B let φz ∈ L2 (∂B) be defined by (6.23); that is, φz (x) = a · ∇z N (x, z), x ∈ ∂B. Then z∈D

⇐⇒

∞  (φz , ψj )2L2 (∂B) j=1

λj

< ∞

(6.28)

or, equivalently, z∈D

⇐⇒

⎤−1 ⎡ ∞  (φz , ψj )2L2 (∂B) ⎦ χ(z) := ⎣ > 0. λj j=1

(6.29)

Here we agreed on the setting that the inverse of the series is zero in the case of divergence. Therefore, χ vanishes outside of D and is positive in the interior of D. The function 1, χ(z) > 0, sign χ(z) = 0, χ(z) = 0, is thus the characteristic function of D. We finish this section with some further remarks. We leave it to the reader to show (see Problems 6.2–6.4) that in the case of B = B(0, 1) being the unit disk and D = B(0, R) the disk of radius R < 1 the ratios (φz , ψj )2L2 (∂B) /λj behave as (|z|/R)2j . Therefore, convergence holds if and only if |z| < R which confirms the assertion of the last theorem. In practice, only finitely many measurements are available; that is, the data operator Λ1 − Λ is replaced by a matrix M ∈ Rm×m . The question of convergence of the series is obsolete because the sum consists of only finitely many terms. However, in practice, it is observed that the value of this sum is much smaller for points z inside of D than for points outside of D. Some authors (see [123]) suggest to test the “convergence”  by determining the slope of the straight line that best fits the curve j → ln (φz , ψj )2L2 (∂B) /λj (for some j only). The points z for which the slope is negative provide a good picture of D. A rigorous justification of a projection method to approximate the (infinite) series by a (finite) sum has been given in [177]. In the implementation of the factorization method, only the relative data operator Λ1 − Λ has to be known and no other information on D. For example, it is allowed (see Assumption 6.10) that D consist of several components. Furthermore, the fact that the medium D is penetrable is not used. If one imposes some boundary condition on ∂D, the same characterization as in Theorem 6.18

6.5

Problems

237

holds. For example, in [123], the factorization method has been justified for an insulating object D. In particular, the factorization method provides a proof of uniqueness of D independent of the nature of D; that is, whether it is finitely conducting, a perfect conductor (Dirichlet boundary conditions on ∂D), a perfect insulator (Neumann boundary conditions on ∂D), or a boundary condition of Robin-type.

6.5

Problems

6.1 Let D be a domain with D ⊂ B and γ ∈ L∞ (B) piecewise constant with γ = γ0 in D for some γ0 ∈ R and γ = 1 in B \ D. Let u ∈ H1 (B) ∩ C 2 (B \ ∂D) be a solution of the variational equation (6.4) and assume that u|D and u|B\D have differentiable extensions to D and B \ D, respectively. Show that u solves Δu = 0 in B \∂D and ∂u/∂ν = f on ∂B and u|+ = u|− on ∂D and γ0 ∂u|− /∂ν = ∂u|+ /∂ν on ∂D. Hint: Use Green’s first theorem. For the following problems let B be the unit disk in R2 with center at the origin. 6.2 Show that the fundamental solution Φ and the Green’s function N are given in polar coordinates (x = r(cos t, sin t) and z = ρ(cos τ, sin τ ) ) as ∞ 1  1 + ρ ,n 1 ln r + cos n(t − τ ) , Φ(x, z) = − 2π 2π n=1 n r  x ,z N (x, z) = Φ(x, z) + Φ |x|2  ∞ 1 1 n 1 1 ln r + ρ = − + rn cos n(t − τ ) , 2π 2π n=1 n rn for ρ = |z| < |x| = r. Hint: Write Φ in the form

  + ρ ,2 1 ρ 1 ln r − ln 1 + − 2 cos(t − τ ) 2π 4π r r   ∞ 1 n 1 2 and show n=1 n α cos(ns) = − 2 ln 1 + α − 2 α cos s by differentiation respect to α and applying the geometric series formula for ∞ with n−1 exp(ins). n=1 α Φ(x, z) = −

6.3 Show that φz from (6.23) is given by φz (x) =

a · (x − z) , π|x − z|2

|x| = 1 .

Also compute φz in polar coordinates for a = (cos α, sin α) by the formulas of Problem 6.2.

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Electrical Impedance Tomography

6.4 Compute the eigenvalues λn and the normalized eigenfunctions ψn ∈ L2 (∂B) of Λ1 − ΛγR and the coefficients (φz , ψn )L2 (∂B) for the case of Example 6.6. Compute the ratios (φz , ψn )2L2 (∂B) /λn and validate the condition (6.29) of Theorem 6.18. 6.5 Consider the case of D ⊂ B being the annulus D = {x ∈ B : R1 < |x| < R2 } for some 0 < R1 < R2 < 1. Compute again the eigenvalues λn and the normalized eigenfunctions ψn ∈ L2 (∂B) of Λ1 − Λ and the coefficients (φz , ψn )L2 (∂B) . Verify that you can only determine the outer boundary {x : |x| = R2 } by the factorization method. 6.6 Let f : [a, b] → R>0 be a Lipschitz continuous function; that is, f (x) > 0 for all x ∈ [a, b] and there exists L > 0 with |f (x) − f (y)| ≤ L|x − y| for all x, y ∈  [a, b]. Define D := {(x1 , x2 ) ∈ [a, b] × R : 0 < x2 < f (x1 ) for x1 ∈ [a, b] . Show that this Lipschitz domain D ⊂ R2 satisfies the exterior cone condition of Assumption 6.10.

Chapter 7

An Inverse Scattering Problem 7.1

Introduction

We consider acoustic waves that travel in a medium such as a fluid. Let v(x, t) be the velocity vector of a particle at x ∈ R3 and time t. Let p(x, t), ρ(x, t), and S(x, t) denote the pressure, density, and specific entropy, respectively, of the fluid. We assume that no exterior forces act on the fluid. Then the movement of the particle is described by the following equations. ∂v 1 + (v · ∇)v + γ v + ∇p ∂t ρ ∂ρ + div (ρv) ∂t f (ρ, S) ∂S + v · ∇S ∂t

=

0

(Euler’s equation) ,

(7.1a)

=

0

(continuity equation) ,

(7.1b)

= p

(equation of state) ,

(7.1c)

=

(adiabatic hypothesis) ,

(7.1d)

0

where the function f depends on the fluid. γ is a damping coefficient, which we assume to be piecewise constant. This system is nonlinear in the unknown functions v, ρ, p, and S. Let the stationary case be described by v0 = 0, timeindependent distributions ρ = ρ0 (x) and S = S0 (x), and constant p0 such that   p0 = f ρ0 (x), S0 (x) . The linearization of this nonlinear system is given by the (directional) derivative of this system at (v0 , p0 , ρ0 , S0 ). For deriving the linearization, we set v(x, t) = ε v1 (x, t) + O(ε2 ) , p(x, t) = p0 + ε p1 (x, t) + O(ε2 ) , ρ(x, t) S(x, t)

= ρ0 (x) + ε ρ1 (x, t) + O(ε2 ) , = S0 (x) + ε S1 (x, t) + O(ε2 ) ,

© Springer Nature Switzerland AG 2021 A. Kirsch, An Introduction to the Mathematical Theory of Inverse Problems, Applied Mathematical Sciences 120, https://doi.org/10.1007/978-3-030-63343-1 7

239

240

Inverse Scattering Problem

and we substitute this into (7.1a), (7.1b), (7.1c), and (7.1d). Ignoring terms with O(ε2 ) leads to the linear system ∂v1 1 + γ v1 + ∇p1 = 0 , ∂t ρ0 ∂ρ1 + div (ρ0 v1 ) = 0 , ∂t ∂f (ρ0 , S0 ) ∂f (ρ0 , S0 ) ρ1 + S 1 = p1 , ∂ρ ∂S ∂S1 + v1 · ∇S0 = 0 . ∂t

(7.2a) (7.2b) (7.2c) (7.2d)

First, we eliminate S1 . Because   ∂f (ρ0 , S0 ) ∂f (ρ0 , S0 ) ∇ρ0 + ∇S0 , 0 = ∇f ρ0 (x), S0 (x) = ∂ρ ∂S we conclude by differentiating (7.2c) with respect to t and using (7.2d)   ∂p1 ∂ρ1 = c(x)2 + v1 · ∇ρ0 , (7.2e) ∂t ∂t where the speed of sound c is defined by c(x)2 :=

 ∂  f ρ0 (x), S0 (x) . ∂ρ

Now we eliminate v1 and ρ1 from the system. This can be achieved by differentiating (7.2e) with respect to time and using equations (7.2a) and (7.2b). This leads to the wave equation for p1 :   1 ∂ 2 p1 (x, t) ∂p1 (x, t) 2 = c(x) ρ0 (x) div ∇p1 (x, t) . + γ (7.3) ∂t2 ∂t ρ0 (x) Now we assume that terms involving ∇ρ0 are negligible and that p1 is timeperiodic; that is, of the form   p1 (x, t) = Re u(x) e−iωt with frequency ω > 0 and a complex-valued function u = u(x) depending only on the spatial variable. Substituting this into the wave equation (7.3) yields the three-dimensional Helmholtz equation for u: ω2  γ Δu(x) + 1+i u = 0. 2 c(x) ω In free space, c = c0 is constant and γ = 0. We define the wave number and the index of refraction by γ c20  ω 1 + i . (7.4) > 0 and n(x) := k := c0 c(x)2 ω

Inverse Scattering Problem

241

The Helmholtz equation then takes the form Δu + k 2 n u = 0

(7.5)

where n is a complex-valued function with Re n(x) ≥ 0 and Im n(x) ≥ 0. This equation holds in every source-free domain in R3 . We assume in this chapter that there exists a > 0 such that c(x) = c0 and γ(x) = 0 for all x with |x| ≥ a; that is, n(x) = 1 for |x| ≥ a. This means that the inhomogeneous medium {x ∈ R3 : n(x) = 1} is bounded and contained in the ball B(0, a) := {y ∈ R3 : |y| < a} of radius a. By B[0, a] := {y ∈ R3 : |y| ≤ a}, we denote its closure. We further assume that the sources lie outside the ball B[0, a]. These sources generate “incident” fields ui , that satisfy the unperturbed Helmholtz equation Δui + k 2 ui = 0 outside the sources. In this introduction, we assume that ui is either a point source or a plane wave; that is, the timedependent incident fields have the form pi1 (x, t) =

1 Re eik|x−z|−iωt ; |x − z|

eik|x−z| , |x − z|

that is, ui (x) =

for a source at z ∈ R3 , or ˆ

pi1 (x, t) = Re eikθ·x−iωt ;

ˆ

that is, ui (x) = eikθ·x ,

for a unit vector θˆ ∈ R3 . In any case, ui is a solution of the Helmholtz equation Δui + k 2 ui = 0 in 3 R \ {z} or R3 , respectively. In the first case, pi1 describes a spherical wave that travels away from the source with velocity c0 . In the second case, pi1 is a plane wave that travels in the direction θˆ with velocity c0 . The incident field is disturbed by the medium described by the index of refraction n and produces a “scattered wave” us . The total field u = ui + us satisfies the Helmholtz equation Δu + k 2 n u = 0 outside the sources; that is, the scattered field us satisfies the inhomogeneous equation Δus + k 2 n us = k 2 (1 − n)ui

(7.6)

where the right-hand side is a function of compact support in B(0, a). Furthermore, we expect the scattered field us to behave as a spherical wave far away from the medium. This can be described by the following radiation condition ∂us (x) − ik us (x) = O(1/r2 ) as r = |x| −→ ∞ , ∂r

(7.7)

uniformly in x/|x| ∈ S 2 . Here we denote by S 2 the unit sphere in R3 . The smoothness of the solution us depends on the smoothness of the refractive index n. We refer to the beginning of Subsection 7.2 for more details. We have now derived a (almost) complete description of the direct scattering problem. Let the wave number k > 0, the index of refraction n ∈ L∞ (R3 ) with n(x) = 1 for |x| ≥ a, and the incident field ui be given. Determine the scattered

242

Inverse Scattering Problem

field us that satisfies the source equation (7.6) in some generalized sense (to make more precise later) and the radiation condition (7.7). In the inverse problem, one tries to determine the index of refraction n from measurements of the field u outside of B(0, a) for several different incident fields ui and/or different wave numbers k. The following example shows that the radially symmetric case reduces to an ordinary differential equation. Example 7.1 Let n = n(r) be radially symmetric: n is independent of the spherical coordinates. Because in spherical polar coordinates (r, φ, θ),



 ∂2 ∂ ∂ 1 1 ∂ 1 ∂ + Δ = 2 r2 + 2 2 sin θ , r ∂r ∂r r2 sin θ ∂θ ∂θ r sin θ ∂φ2 the Helmholtz equation for radially symmetric u = u(r) reduces to the following ordinary differential equation of second order, 1  2   r u (r) + k 2 n(r) u(r) = 0 ; r2 that is, 2  u (r) + k 2 n(r) u(r) = 0 for r > 0 . (7.8a) r From the theory of linear ordinary differential equations of second order with singular coefficients, we know that in a neighborhood of r = 0 there exist two linearly independent solutions, a regular one and one with a singularity at r = 0. We construct them by making the substitution u(r) = v(r)/r in (7.8a). This yields the equation u (r) +

v  (r) + k 2 n(r) v(r) = 0

for r > 0 .

(7.8b)

For the simplest case, where n(r) = 1, we readily see that u1 (r) = α sin(kr)/r and u2 (r) = β cos(kr)/r are two linearly independent solutions. u1 is regular and u2 is singular at the origin. Neither of them satisfies the radiation condition. However, the combination u(r) = γ exp(ikr)/r does satisfy the radiation condition because

 1 exp(ikr)  u (r) − iku(r) = −γ = O 2 r2 r as is readily seen. For the case of arbitrary n, we construct a fundamental system {v1 , v2 } of (7.8b) (compare with Section 5.2); that is, v1 and v2 satisfy (7.8b) with v1 (0) = 0, v1 (0) = 1, and v2 (0) = 1, v2 (0) = 0. Then u1 (r) = v1 (r)/r is the regular and u2 (r) = v2 (r)/r is the singular solution. In the next section, we rigorously formulate the direct scattering problem and prove the uniqueness and existence of a solution. The basic ingredients for the uniqueness proof are a result by Rellich (see [222]) and a unique continuation principle for solutions of the Helmholtz equation. We prove neither

7.2

The Direct Scattering Problem

243

Rellich’s lemma nor the general continuation principle, but rather give a simple proof for a special case of a unique continuation principle that is sufficient for the uniqueness proof of the direct problem. This suggestion was made by H¨ahner (see [116]). We then show the equivalence of the scattering problem with an integral equation. Existence is then proven by an application of the Riesz theorem A.36 of Appendix A. Section 7.3 is devoted to the introduction of the far field patterns that describe the scattered fields “far away” from the medium. We collect some results on the far field operator, several of which are needed in Sections 7.5 and 7.7. The question of injectivity of the far field operator is closely related to an unusual eigenvalue problem which we call the interior transmission eigenvalue problem. We will investigate this eigenvalue problem in Section 7.6. In Section 7.4, we prove uniqueness of the inverse problem. Section 7.5 is devoted to the factorization method which corresponds to the method in Section 6.4 and provides a very simple characterization of the support of the contrast by the far field patterns. This method is rigorously justified under the assumption that, again, the wavenumber is not an interior transmission eigenvalue. Since the interior transmission eigenvalue problem is also an interesting problem in itself and widely studied during the past decade we include Section 7.6 for some aspects of this eigenvalue problem. Finally, in Section 7.7, we present three classical numerical algorithms for solving the inverse scattering problem.

7.2

The Direct Scattering Problem

In this section, we collect properties of solutions to the Helmholtz equation that are needed later. We prove uniqueness and existence of the direct scattering problem and introduce the far field pattern. In the remaining part of this chapter, we restrict ourselves to scattering problems for plane incident fields. Throughout this chapter, we make the following assumptions. Let n ∈ L∞ (R3 ) and a > 0 with n(x) = 1 for almost all |x| ≥ a. Assume that Re n(x) ≥ 0 and Im n(x) ≥ 0 for almost all x ∈ R3 . Let k ∈ R, k > 0, ˆ = 1. We set ui (x) := exp(ik θˆ · x) for x ∈ R3 . Then ui solves and θˆ ∈ R3 with |θ| the Helmholtz equation Δui + k 2 ui = 0

in R3 .

(7.9)

We again formulate the direct scattering problem. Given n, k, θˆ satisfying the 2 previous assumptions, determine u ∈ Hloc (R3 ) such that Δu + k 2 n u = 0

in R3 ,

(7.10)

and us := u − ui satisfies the Sommerfeld radiation condition ∂us − ikus = O(1/r2 ) for r = |x| → ∞ , ∂r

(7.11)

uniformly in x/|x| ∈ S 2 . Since the index function n is not smooth we cannot expect that the solution u is smooth either. Rather, it belongs to the (local)

244

Inverse Scattering Problem

2 Sobolev space Hloc (R3 ). We recall that for any open set Ω ⊂ R3 and p ∈ N the p Sobolev space H (Ω) is defined as the completion of

u ∈ C p (Ω) :

∂ |j|1

u ∈ L2 (Ω) for all j ∈ N30 with |j|1 ≤ p ∂xj11 ∂xj22 ∂xj33

with respect to the norm

u H p (Ω)

       ∂ |j|1 u(x) 2     =   ∂xj1 ∂xj2 ∂xj3  dx .  1 2 3  j ∈ N30 Ω |j|1 ≤ p

Here we have set |j|1 = |j1 | + |j2 | + |j3 | for j = (j1 , j2 , j3 ) ∈ N30 . We refer to Chapter 6 where we already used Sobolev spaces of functions on two-dimensional p (Ω) are defined by domains B. The local spaces Hloc   p Hloc (Ω)= u : Ω→C : u|B ∈ H p (B) for every bounded domain B with B ⊂ Ω . For us, the spaces H 1 (Ω) and H 2 (Ω) (and their local analogies) are particularly important. We define the subspace H0p (Ω) of H p (Ω) by the closure of the set C0p (Ω) = ψ ∈ C p (Ω) : ψ has compact support in Ω in H p (Ω). By definition it is closed, and one can show that it is a strict subspace of H p (Ω). The space H01 (Ω) models the class of functions which vanish on the boundary of Ω while for functions in H02 (Ω) also the normal derivatives ∂ψ/∂ν vanish at ∂Ω. The following versions of Green’s theorem are not difficult to prove by approximating u and v by sequences of smooth functions. Lemma 7.2 Let Ω ⊂ R3 be a bounded domain. Then    u Δv + ∇u · ∇v dx = 0 for all u ∈ H01 (Ω), v ∈ H 2 (Ω) ,

(7.12a)

Ω



  u Δv − v Δu dx = 0

for all u ∈ H 2 (Ω), v ∈ H02 (Ω) .

(7.12b)

Ω

The application of (7.12a) to solutions of the Helmholtz equation yields the following version.   Lemma 7.3 Let v, w ∈ H 2 B(0, b) ) be solutions of the Helmholtz equation Δu + k 2 u = 0 in some annular region A = {x ∈ R3 : a < |x| < b}. Then v and w are analytic in A, and for every R ∈ (a, b) it holds that     ∂w ds = ∇v · ∇w + v Δw dx . (7.13) v ∂r |x|=R

|x| a, and the radiation condition (7.11) is well-defined. Proof: The smoothness of v and w follow from general regularity results for solutions of the Helmholtz equation Δu + k 2 u = 0 and is not proven here, see, e.g., [161], Corollary 3.4, or [55], Theorem 2.2. To show (7.13) we choose support in B(0, b) such that ρ(x) = 1 for |x| ≤ R. ρ ∈ C ∞ (R3 ) with  compact  Then ρv ∈ H01 B(0, b) as easily seen and thus by (7.12a)    0 = ρv Δw + ∇w · ∇(ρv) dx |x| 0 we estimate   −|j|2 + i j · (2tˆ e − ip) − (it + α)      1 ≥ Re[· · · ] = |j|2 + 2tj2 − j · p + α ≥ 1 + |j|2 2 for all j ∈ Z3 with |j| ≥ j0 for some j0 ∈ N. Furthermore,   −|j|2 + i j · (2tˆ e − ip) − (it + α) ≥ | Im[· · · ]| = t|2j1 − 1| ≥ t for all j ∈ Z3 and t > 0. Therefore, wj =

gj −|j|2 + i j · (ip + 2tˆ e) − (it + α)

2 (Q) because are well defined for all j ∈ Z3 and w ∈ Hper   [1 + |j|2 ]2 |wj |2 ≤ 4 |gj |2 . |j|≥j0

|j|≥j0

Furthermore, the solution operator  gj ei j·x , (Lt g)(x) := 2 + i j · (ip + 2tˆ −|j| e ) − (it + α) 3

g ∈ L2 (Q),

j∈Z

is bounded from L2 (Q) into itself with Lt L(L2 (Q)) ≤ 1/t for every t > 0.  Now we can give a simple proof of the following version of a unique continuation principle. Theorem 7.7 Let n ∈ L∞ (R3 ) with n(x) = 1 for |x| ≥ a be given. Let u ∈ H 2 (R3 ) be a solution of the Helmholtz equation Δu + k 2 n u = 0 in R3 such that u(x) = 0 for all |x| ≥ b for some b ≥ a. Then u has to vanish in all of R3 . Proof: Define eˆ = (1, i, 0) ∈ C3 as before, set ρ = 2b/π, and define the function w(x) := ei/2 x1 −t eˆ·x u(ρx), x ∈ Q := (−π, π)3 ,

248

Inverse Scattering Problem

for some t > 0. Then w(x) = 0 for all |x| ≥ π/2, in particular near the boundary of the cube Q. Extend w to a 2π-periodic function in R3 by w(2πj + x) := w(x) 2 (Q), and w satisfies the for x ∈ Q and all j ∈ Z3 , j = 0. Then w ∈ Hper differential equation Δw + (2tˆ e − ip) · ∇w − (it + 1/4) w = −ρ2 k 2 n ˜w. Here, we have set p = (1, 0, 0) and n ˜ (2πj + x) := n(ρx) for almost all x ∈ [−π, π]3 and j ∈ Z3 . Application of the previous lemma to this differential equation yields the existence of a linear bounded operator Lt from L2 (Q) into itself with Lt L(L2 (Q)) ≤ 1/t such that the differential equation is equivalent to   ˜w . w = −ρ2 k 2 Lt n Estimating

w L2 (Q) ≤

ρ2 k 2 ρ2 k 2 n ∞

˜ nw L2 (Q) ≤

w L2 (Q) t t

yields w = 0 for sufficiently large t > 0. Thus, also u has to vanish.



The preceding theorem is a special case of a far more general unique continuation principle, which we formulate without proof here. 2 (Ω) be a solution of the Helmholtz equation Δu + k 2 nu = 0 in a Let u ∈ Hloc 3 domain Ω ⊂ R (i.e., Ω is open and connected). Furthermore, let n ∈ L∞ (Ω) and u(x) = 0 on some open set. Then u = 0 in all of Ω. For a proof we refer to, for example, [55]. Now we can prove the following uniqueness result. Theorem 7.8 (Uniqueness) The problem (7.10), (7.11) has at most one solution; that is, if u is a solution corresponding to ui = 0, then u = 0. Proof:

Let ui = 0. The radiation condition (7.11) yields 2

O(1/R )

 = |x|=R



= |x|=R

2    ∂u    ∂r − ik u ds

(7.18)

 2    ∂u    + k 2 |u|2 ds + 2 k Im  ∂r 

|x|=R

u

∂u ds. ∂r

We transform the last integral using Green’s formula (7.13) for v = u and w = u; that is,     ∂u |∇u|2 − k 2 n |u|2 dx , ds = u ∂r |x|=R

|x| 0 (depending  only on n, D, and A) such that w H (A) ≤  c f L2 (D) + w L2 (D) .

(b) Every solution w ∈ H 2 (D) of Δw + k 2 nw = f in D is also an ultra-weak solution. Proof: (a) Let A be any open bounded set such that A ⊂ D. We choose ρ ∈ C ∞ (D) with compact support in D such that ρ = 1 on A. Furthermore, let Q be a cube containing D in its interior. For any φ ∈ C ∞ (R3 ) we take ψ = ρφ in (7.21). With Δψ = ρΔφ + 2 ∇ρ · ∇φ + φΔρ we have    ρ w φ − Δφ dx D

 = D



=

  2   w (k n + 1)ρφ + 2 ∇ρ · ∇φ + φΔρ − f ρ φ dx 

(7.22)

 g φ + h · ∇φ dx

D

with g = w[(k 2 n + 1)ρ + Δρ] − f ρ and h = 2w∇ρ. We can replace the region of integration by Q because ρ vanishes in D \ Q. Without loss of generality we L2 (Q) we expand them into Fourier assume that Q= (−π, π)3 . Since g, h ∈  ij·x and h(x) = j∈Z3 hj eij·x with gj ∈ C and hj ∈ series g(x) = j∈Z3 gj e   3 2 2 C such that < ∞ and 3 |gj | j∈Z j∈Z3 |hj | < ∞ and make the ansatz  (ρw)(x) = j∈Z3 wj eij·x in Q. We take φ(x) = e−i·x for some  ∈ Z3 in (7.22) and have (1 + ||2 ) w = g − i  · h and thus (1 + ||2 ) |w |2 ≤

 2    2 |g | + ||2 |h |2 ≤ 2 |g |2 + |h |2 2 1 + ||

(7.23)

1 Therefore, ρw ∈ Hper (Q) and

w H 1 (A)

≤ ≤

  1 (Q) ≤ c2 g L2 (Q) + h L2 (Q) c1 ρw Hper   c3 w L2 (D) + f L2 (D) .

1 Since A was arbitrary we have shown that w ∈ Hloc (D). Now we repeat the first part of the proof but apply Green’s first formula to the second term of the

7.2

The Direct Scattering Problem

251

right hand side of (7.22). This yields       ρ w φ − Δφ dx = φ w(k 2 n + 1)ρ − 2 div (w∇ρ) + w Δρ − f ρ dx .    Q

=: g

Q

Now we argue in the same way but with h = 0. Estimate (7.23) yields (1 + 2 ||2 )2 |w |2 ≤ 2|g |2 and thus ρw ∈ Hper (Q) and ≤

w H 2 (A)

≤

2 (Q) ≤ c2 g L2 (Q) c1 ρw Hper   c3 w L2 (D) + ∇w L2 (B) + f L2 (D)

where B is the support of ρ. Now we substitute the estimate for w H 1 (B) which yields the desired estimate. (b) This follows directly from Green’s theorem in the form (7.12b).  Now we construct volume potentials with the fundamental solution (7.19). Theorem 7.11 Let Ω ⊂ R3 be a bounded domain. For every φ ∈ L2 (Ω) the volume potential  φ(y) Φ(x, y) dy , x ∈ R3 , (7.24) v(x) := Ω 2 Hloc (R3 )

yields a function v ∈ that satisfies the radiation condition (7.11) and is the only radiating2 solution of Δv + k 2 v = −φ. Furthermore, for every ball B = B(0, R) containing Ω in its interior there exists c > 0 (only dependent on B, k, and Ω) such that (7.25)

v H 2 (B) ≤ c φ L2 (Ω) . Proof: First we state without proof (see, e.g., [161], Theorem 3.9) that for any k ∈ C and φ ∈ C01 (Ω) ; that is, φ ∈ C 1 (Ω) with compact support in Ω, the potential v is in C 2 (R3 ) and solves Δv + k 2 v = −φ in Ω and Δv + k 2 v = 0 in the exterior of Ω. Second, we fix φ ∈ L2 (Ω) and choose a sequence φj ∈ C01 (Ω) which converges to φ in L2 (Ω). Let v and vj be the corresponding potentials, and let ψ ∈ C ∞ (R3 ) some test function with compact support. Then Δvj + k 2 vj = −φj in R3 and thus by Green’s second theorem   vj (Δψ + k 2 ψ) dx = − φj ψ dx . R3

Ω

Let B(0, R) be a ball that contains the support of ψ. From  the boundedness of the volume integral operator from L2 (Ω) into L2 B(0, R) we conclude that vj converges to v in L2 (B(0, R)) as j tends to infinity. Therefore,   v (Δψ + k 2 ψ) dx = − φ ψ dx R3 2 that

is, it satisfies the radiation condition (7.11)

Ω

252

Inverse Scattering Problem

for all ψ ∈ C ∞ (R3 ) with compact support. Therefore, v is an ultra-weak solution of Δv + k 2 v = −φ in R3 . The  regularity result of Lemma 7.10 applied to 2 B(0, R + 1)) and the estimate D = B(0, R + 1) yields v ∈ Hloc  

v H 2 (B(0,R) ≤ c1 v L2 (B(0,R+1) + φ L2 (Ω) ≤ c2 φ L2 (Ω) 2 where  we used again the boundedness of the volume potential from L (Ω) to 2  L B(0, R + 1) .

Now we can transform the scattering problem into a Fredholm integral equation of the second kind. The following theorem is needed quite often later on. 2 Theorem 7.12 (a) Let u ∈ Hloc (R3 ) be a solution  of the scattering problem (7.10), (7.11). Then u|B(0,a) belongs to L2 B(0, a) and solves the Lippmann– Schwinger integral equation    i 2 u(x) = u (x) − k 1 − n(y) Φ(x, y) u(y) dy , x ∈ B(0, a) . (7.26) |y| 0. Lemma 7.44 The characteristic function d from (7.87) is an even entire func1 tion of order one provided n(1) = 1 or η := 0 n(s)ds = 1. Proof: d is an even function because yk and also k → sin k/k and cos are all functions of k 2 . Furthermore, from the Volterra integral equation (5.12b) for the part u2 (·, λ, q) of the fundamental system {u1 , u2 } of (5.7a), (5.7b) we derive easily that λ → u2 (η, λ, q) is holomorphic in all of C and thus also k → yk (1) = n(1)−1/4 zk (η) = [n(0)n(1)]−1/4 u2 (η, k 2 , q). The same holds for the derivative. Therefore, d is an entire function. It remains to compute the order of d. From (7.90) and the estimate | sin z| ≤ e|z| we observe that |d(k)| ≤ c e|k|(η+1) for some c > 0. Therefore, for |k| = r, ln[ln c + (η + 1)r] ln[(η + 2)r] ln(η + 2) ln(ln |d(k)|) ≤ ≤ = + 1 ln r ln r ln r ln r for large values of r which proves that ρ ≤ 1. Now we set k = ti for t ∈ R>0 . t(η+1) If n(1) = 1 then it is easy to see that |f (it)| ≥ |A| for large values of t 2 e (where f if given below (7.90)) and thus, again from (7.90), |d(it)| ≥ c et(η+1) for some c > 0. This yields that ρ ≥ 1 and ends the proof if n(1) = 1. If n(1) = 1 then η = 1 and A = 0 and one gets |d(it)| ≥ c et|η−1| and again ρ ≥ 1.  The following theorem is a special case of the more general factorization theorem of Hadamard which we cite without proof (see, e.g., [19], Chapter 2).

288

Inverse Scattering Problem

Theorem 7.45 (Hadamard) Let f be an entire function of order one, let m ≥ 0 be the order7 of the root z = 0 of f (in particular, m = 0 if f (0) = 0), and let {aj : j ∈ N} be the nonzero roots of f repeated according to multiplicity. Then there exists a polynomial p of degree at most one such that f (z) = ep(z) z m

∞

% j=1

1−

z aj

 ez/aj ,

z ∈ C.

(7.92)

With this theorem we can prove a slight generalization of a theorem of Laguerre. Theorem 7.46 Let f be an entire function of order one which is real for real values of z. Suppose that f has infinitely many real zeros and only a finite number of complex ones. Then f has a single critical point8 on each interval formed by two consecutive real zeros of f provided this interval is sufficiently far away from the origin. Proof: Since f has only finitely many complex zeros (which will occur in conjugate pairs) there exists a polynomial q whose roots are exactly all those complex ones of f as well as a possible root at z = 0 which has the same order as the possible zero of f at z = 0. We factorize f (z) = q(z)g(z) where all the roots of g are given by {aj : j ∈ N} in nondecreasing magnitude (where multiple roots are repeated according to their multiplicities). We apply Hadamard’s theorem to the function g and have for real values of x ∈ / {aj : j ∈ N} g  (x) g(x)

  

∞  d d x x  = ln g(x) = p (x) + + ln 1 − dx dx aj aj j=1   ∞  1 1 = α + + x − a a j j j=1

where α = p (x) is a real constant (since p is a polynomial of order at most one). Differentiating this expression yields





 d f  (x) d q  (x) d g  (x) = + dx f (x) dx q(x) dx g(x)

  ∞  1 d q (x) , x∈ / {aj : j ∈ N} . = − dx q(x) (x − aj )2 j=1 Since q is a polynomial   with no real zeros (except possibly z = 0) there exists  d q (x)  c > 0 with  dx q(x)  ≤ |x|c 2 for all |x| ≥ 1. Choose N ∈ N with N > 2c and 7 Note 8 that

that the order of a root of a holomorphic function is always finite. is, a zero of f 

7.6

The Interior Transmission Eigenvalue Problem

289

then R > 0 with |aj − x|2 ≤ 2|x|2 for all |x| ≥ R and j = 1, . . . , N . Then d dx



f  (x) f (x)

 ≤

N  1 c N 2c − N c − ≤ − = < 0 2 2 2 2 |x| (x − aj ) |x| 2|x| 2|x|2 j=1

for |x| ≥ R, x ∈ / {aj : j ∈ N}. Therefore, f  /f is strictly decreasing in every interval (a , a+1 ) with a < a+1 and  large enough. Furthermore, from f  (x) f (x)

 ∞   1 q  (x) g  (x) q  (x) 1 = + = + α + + q(x) g(x) q(x) x − aj aj j=1    1 1 q  (x) + α + + = q(x) x − aj aj j =,+1     1 1 1 1 + + + + x − a a x − a+1 a+1 f  (x) x→a + f (x)

we conclude that lim

= +∞ and

 lim f (x) x→a+1 − f (x)

(7.93)

= −∞. Therefore, f  /f

has exactly one zero in the interval (a , a+1 ) which ends the proof.



As a corollary we obtain Laguerre’s theorem (see [19], Chapter 2). Corollary 7.47 (Laguerre) Let f be an entire function of order one which is real for real values of z and all of its zeros are real. Then all of the zeros of f  are real as well and interlace those of f . Proof: In this case we set q(z) = z m where m is the order of the zero z = 0. Then, by analytic extension, (7.93) holds also for complex z instead of x ∈ R. If z ∈ C is a critical point then from (7.93)   ∞  ∞    1 z − aj 1 1 m mz + α + + + α + + 0 = = . z z − aj aj |z|2 |z − aj |2 aj j=1 j=1 Taking the imaginary part yields 0 = − Im z

∞  1 m − Im z |z|2 |z − aj |2 j=1

which implies that z is real. Now we follow exactly the proof of the preceding theorem.  Now we are able to answer the question of existence of complex eigenvalues for the special case of n being constant. For constant n the function yk from

290

Inverse Scattering Problem

(7.86) is given by yk (r) = 1 kη fη (k) where

sin(kηr) kη

for r ≥ 0 where η =

fη (k) = η cos(kη) sin k − cos k sin(kη) ,

√

n. Therefore, d(k) = k ∈ C.

(7.94)

Here we indicated the dependence on η > 0. For the author the following result is rather surprising. √ Theorem 7.48 If η = n = 1 is an integer or the reciprocal of an integer then no complex eigenvalues with eigenfunctions depending only on r exist. Otherwise, there exist infinitely many complex eigenvalues. Proof: We consider the function fη from (7.94) and note first that −ηf1/η (ηk) = fη (k) for all k ∈ C. Therefore, a complex zero of fη exists if, and only if, a complex zero of f1/η exists. It is thus sufficient to study the case η > 1. Let first η > 1 be an integer. Then fη (k) = 0 if, and only if, k is a critical 2  point of the entire function g(k) = sin(ηk) sin k because g (k) = fη (k)/ sin k. Since all of the zeros of g are real, by Corollary 7.47 also its critical points are real; that is, all of the zeros of fη are real. Let now η > 1 not be an integer. We will construct a sequence I of intervals which tend to infinity such that fη does not change its sign on I and each I contains two consecutive real critical points of fη . By Theorem 7.46 this is only possible if there exist infinitely many complex zeros of fη . From fη (k) = (1 − η 2 ) sin(kη) sin k  jπ  we observe that η : j ∈ N and {jπ : j ∈ N} are the critical points of / N for all fη . Choose a sequence m ∈ N converging to infinity such that η m ∈  ∈ N. Fix  ∈ N in the following and set m = m for abbreviation. The interval (mη − 1, mη) contains an integer j. Set ε = mη − j ∈ (0, 1). The two points jπ η π π and mπ are consecutive zeros of fη because their distance is mπ − jπ η = εη < η. Furthermore,

 jπ jπ jπ − cos sin(jπ) fη = η cos(jπ) sin η η η

 επ επ and = η (−1)j sin mπ − = η (−1)j+1+m sin η η fη (mπ) = η cos(mπη) sin(mπ) − cos(mπ) sin(mπη) = (−1)m+1 sin(jπ + επ) = (−1)m+1+j sin(επ) .  and fη (mπ) coincide because ε, ε/η ∈ We observe that the signs of fη jπ η (0, 1). Furthermore, fη has no zero between jπ η and mπ because otherwise there would be another critical point between them. Therefore, the interval  I = jπ η , mπ has the desired properties. We refer to Problem 7.5 for two explicit examples where the assertions of this theorem are illustrated.

7.6

The Interior Transmission Eigenvalue Problem

291

There exist also results on the existence of complex transmission eigenvalues for non-constant refractive indices under  certain conditions.For example, as shown in [55] (see also [58]), if either 1 < n(1) < η or η < n(1) < 1 (where η is again given by (7.88d)) there exist infinitely many real and infinitely many complex eigenvalues. Also, all complex eigenvalues lie in a strip around the real axis if n(1) = 1. For a more detailed investigation of the location for the constant case we refer to [257] and to [55] and the references therein.

7.6.2

Discreteness And Existence in the General Case

We continue now with the general case; that is, D is a bounded Lipschitz domain and n is real-valued with n(x) ≥ 1 + q0 on D for some q0 > 0. For the definition of a Lipschitz domain we refer again to [191] or [161]. For Lipschitz domains we can give an alternative characterization of an interior transmission eigenvalue. We allow k = 0 to be complex. Theorem 7.49 Let D be a Lipschitz domain. k ∈ C \ {0} is an interior transmission eigenvalue if and only if there exist u ∈ H02 (D) and v ∈ L2 (D), not vanishing simultaneously, such that Δu + k 2 nu = k 2 (n − 1) v a.e. in D

and

Δv + k 2 v = 0 in D in the ultra weak sense; that is,    v Δϕ + k 2 ϕ dx = 0 for all ϕ ∈ H02 (D) .

(7.95a)

(7.95b)

D

Proof: Let first u ∈ H02 (D) and v ∈ L2 (D) with (7.95a), (7.95b). Set w = v − u and let φ, ψ ∈ H 2 (D) with φ − ψ ∈ H02 (D). Then 

    v Δφ + k 2 φ − w Δψ + k 2 nψ dx

D

 =

  v Δ(φ − ψ) + k 2 (φ − ψ) dx +

D



  k 2 v (1 − n) ψ + u Δψ + k 2 nψ dx

D

The first integral on the right hand side vanishes because v is an ultra weak solution of Δv +k 2 v = 0. For the second integral we use Green’s second formula (7.12b) which yields  D

  u Δψ + k 2 nψ dx =

 D

  ψ Δu + k 2 nu dx = k 2

 ψ (n − 1) v dx . D

Therefore, the pair (v, w) ∈ L2 (D) × L2 (D) satisfies (7.84). This proves the first direction.

292

Inverse Scattering Problem

For the reverse direction let v, w ∈ L2 (D) with (7.84). Define u = v − w in D and u = 0 in R3 \ D. Furthermore, set ϕ = ψ for some ψ ∈ C ∞ (R3 ) with compact support in (7.84). Then     u Δψ + k 2 nψ dx = k 2 (n − 1) v ψ dx ; R3

D

that is, u is an ultra weak solution of Δu + k 2 nu = k 2 (n − 1)v in R3 . The regularity result of Lemma 7.10 yields u ∈ H 2 (R2 ) with u = 0 in R3 \ D; that is, u ∈ H02 (D) by Lemma 7.4 and Δu + k 2 nu = k 2 (n − 1)v almost  everywhere  in D. Finally, let ϕ ∈ H02 (D) and set ψ = 0 in (7.84). Then D v Δϕ + k 2 ϕ dx = 0. This ends the proof.  This theorem makes it possible to eliminate the function v from the system. Indeed, let u ∈ H02 (D) and v ∈ L2 (D) satisfy (7.95a) and (7.95b). We devide (7.95a) by n − 1 (note that n(x) − 1 ≥ q0 on D by assumption), multiply the resulting equation by Δψ + k 2 ψ for some ψ ∈ H02 (D) and integrate. This gives      dx   = k 2 v Δψ + k 2 ψ dx = 0 Δu + k 2 nu Δψ + k 2 ψ n−1 D

D

by (7.95b); that is, replacing ψ by its complex conjugate,     dx = 0 for all ψ ∈ H02 (D) . Δu + k 2 nu Δψ + k 2 ψ n−1

(7.96)

D

1 On the other side, if u ∈ H02 (D) satisfies (7.96) then we set v = k2 (n−1) [Δu + 2 2 k nu] ∈ L (D) which satisfies (7.95b). In this sense the system (7.95a), (7.95b) is equivalent (7.96) is the weak form of the fourth order  1  Equation   to (7.96). Δu + k 2 nu = 0. equation Δ + k 2 n−1 With respect to (7.96) it is convenient to introduce a new inner product (·, ·)∗ in H02 (D) by  dx , u, ψ ∈ H02 (D) . (u, ψ)∗ := Δu Δψ (7.97) n−1 D

Lemma 7.50 (·, ·)∗ defines an inner product in H02 (D) with corresponding norm

· ∗ which is equivalent to the ordinary norm · H 2 (D) in H02 (D). Proof: By the definition of H02 (D) and a denseness argument it is sufficient to prove the existence of c1 , c2 > 0 with c1 ψ ∗ ≤ ψ H 2 (D) ≤ c2 ψ ∗ for all ψ ∈ C0∞ (D). The left estimate is obvious because n − 1 ≥ q0 > 0 on D. To prove the right estimate we choose a cube Q ⊂ R3 which contains D in its interior, extend any ψ ∈ C0∞ (D) by zero into Q, and then periodically (with respect to Q) into R3 . Without loss of generality we assume that Q = (−π, π)3 .

7.6

The Interior Transmission Eigenvalue Problem

293

2 2 (Q) where c > 0 is independent of Then ψ ∈ Hper (Q) and ψ H 2 (D) ≤ c ψ Hper ψ. With the Fourier coefficients ψj of ψ (see (7.15b)) we have  = [1 + |j|2 ]2 |ψj |2

ψ 2Hper 2 (Q) j∈Z3



= |ψ0 |2 +

[1 + 2|j|2 + |j|4 ]|ψj |2

0 =j∈Z3

≤

|ψ0 |2 + 4



|j|4 |ψj |2 = |ψ0 |2 +

0 =j∈Z3

≤

|ψ0 |2 + 4

4

Δψ 2L2 (Q) (2π)3

n − 1 ∞

ψ 2∗ . (2π)3

 ∈ ∂Q we conclude from the It remains to estimate |ψ 0 |. For x = (π, 0, 0) ij1 π = 0 and thus boundary condition that j ψj e



|ψ0 | ≤

|ψj | =

&

j =0

0 =j∈Z3

as before.

 1 |j|2 |ψj | ≤ |j|2

 1 & |j|4 |ψj |2 ≤ c ψ ∗ |j|4 j =0

j =0



Now we rewrite (7.96) in the form (u, ψ)∗ + k 2 a(u, ψ) + k 4 b(u, ψ) = 0 for all ψ ∈ H02 (D) , where

 a(u, ψ)

=

[n u Δψ + ψ Δu] D



= D

=

nuψ D

dx n−1

dx + [u Δψ + ψ Δu] n−1



b(u, ψ)

(7.98)

 u Δψ dx , D

dx n−1

for u, ψ ∈ H02 (D). The sesqui-linear forms a and b are hermetian. For a this is seen from the second form and Green’s second identity (7.12b). By the representation theorem of Riesz (Theorem A.23 of the Appendix) in the Hilbert space H02 (D) for every u ∈ H02 (D) there exists a unique u ∈ H02 (D) with a(u, ψ) = (u , ψ)∗ for all ψ ∈ H02 (D). We define the operator A from H02 (D) into itself by Au = u . Therefore, a(u, ψ) = (Au, ψ)∗ for all u, ψ ∈ H02 (D). Then it is not difficult to prove that A is linear, self-adjoint, and compact.9 9 For the proof of compactness one needs, however, the fact that H 1 (D) is compactly 0 imbedded in L2 (D), see Problem 7.1.

294

Inverse Scattering Problem

Analogously, there exists a linear, self-adjoint, and compact operator B from H02 (D) into itself with b(u, ψ) = (Bu, ψ)∗ for all u, ψ ∈ H02 (D). Then (7.96) can be written as (u, ψ)∗ + k 2 (Au, ψ)∗ + k 4 (Bu, ψ)∗ = 0 for all ψ ∈ H02 (D); that is, u + k 2 Au + k 4 Bu = 0

in H02 (D) .

(7.99)

This is a quadratic eigenvalue problem in the parameter τ = k 2 . We can reduce it to a linear eigenvalue problem for a compact operator. Indeed, since B is positive (that is, (Bu, u)∗ > 0 for all u = 0) there exists a positive and compact operator B 1/2 from H02 (D) into itself with B 1/2 B 1/2 = B (see formula (A.47) of the Appendix). If u satisfies (7.99) for some k ∈ C \ {0} then set u1 = u and  u2 = k 2 B 1/2 u. Then the pair uu12 ∈ H02 (D) × H02 (D) satisfies

 u1 A + k2 −B 1/2 u2

B 1/2 0



u1 u2

 =

 0 . 0

(7.100)

Therefore 1/k 2 is an eigenvalue of the compact (but not self-adjoint) operator

  −A −B 1/2 . Conversely, if uu12 ∈ H02 (D) × H02 (D) satisfies (7.100) 1/2 B 0 then u = u1 satisfies (7.99). Therefore, well-known results on the spectrum of compact operators (see, e.g., [151]) imply the following theorem. Theorem 7.51 There exists at most a countable number of eigenvalues k ∈ C with no accumulation point in C. By different methods the discreteness can be shown under the weaker assumption that n > 1 only in a neighborhood of the boundary ∂D (see, e.g., [256, 159] or [34]). The question of the existence of real eigenvalue was open for about 20 years.10 . The idea of the proof of the following result goes back to P¨ aiv¨ arinta and Sylvester ([213]) and was generalized to general refractive indices by Cakoni, Gintides, and Haddar ([36]). Theorem 7.52 There exists a countable number of real eigenvalues k > 0 which converge to infinity. Proof: With the Riesz representations A and B of the hermetian forms a and b, respectively, from above we define the family of operators Φ(κ) = I + κA + κ2 B for κ ≥ 0 with corresponding sesqui-linear forms     dx Δu + κnu Δψ + κψ φ(κ; u, ψ) = n−1 D      dx   − κ u Δψ + κψ dx ; (7.101) Δu + κu Δψ + κψ = n−1 D

10 Note

D

that the matrix operator in (7.100) fails to be self-adjoint!

7.6

The Interior Transmission Eigenvalue Problem

295

  that is, φ(κ; u, ψ) = Φ(κ)u, ψ ∗ for all u, ψ ∈ H02 (D). Here, κ plays the role of k 2 and is considered as a parameter. We search for parameters κ > 0 such that 0 is an eigenvalue of Φ(κ); that is, −1 is an eigenvalue of the compact and self-adjoint operator κA + κ2 B. For κ ≥ 0 let λj (κ) ∈ R, j ∈ N, be the nonzero eigenvalues of the compact and self-adjoint operators κA + κ2 B. They converge to zero as j tends to infinity (if there exist infinitely many). The corresponding eigenspaces are finite− dimensional. Let the negative eigenvalues be sorted as λ− 1 (κ) ≤ λ2 (κ) ≤ · · · < 0 where the entries are repeated with its multiplicity. In general, they could be none or finitely many or infinitely many of them. Let m ∈ N be any natural number. First we construct κ ˆ > 0 and a subspace κ; u, u) ≤ 0 for all u ∈ Vm : We choose Vm of H02 (D) of dimension m such that φ(ˆ ε > 0and m pairwise disjoint balls Bj = B(zj , ε), j = 1, . . . , m, of radius ε m with j=1 Bj ⊂ D. Setting n0 := 1 + q0 we note that n(x) ≥ n0 on D. In every ball Bj we consider the interior transmission eigenvalue problem with constant refractive index n0 . By Theorem 7.41 infinitely many real and positive transmission eigenvalues exist. Note that these eigenvalues do not depend on j because a translation of a domain results in the same interior transmission eigenvalues. Let kˆ > 0 be the smallest one with corresponding eigenfunctions uj ∈ H02 (Bj ). Set κ ˆ = kˆ2 and let φj (κ; u, ψ) be the sesqui-linear form corresponding to n0 in κ; uj , ψ) = 0 for all ψ ∈ H02 (Bj ). We extend each uj by zero to a Bj . Then φj (ˆ function in H02 (D). (We use Lemma 7.4 again.) Then {uj : j = 1, . . . , m} are certainly linearly independent because their supports are pairwise disjoint. We : j = 1, . . . , m} as a m−dimensional subspace of H02 (D) define Vm = span{uj  m and compute for u = j=1 αj uj ∈ Vm , using (7.101), φ(ˆ κ; u, u)

=

=

m  j=1 m 

|αj |2 φ(ˆ κ; u j , uj ) |αj |2

j=1

≤

m  j=1

=

m 

 Bj

|αj |

2



Bj

 2 dx Δuj + κ − κ ˆ ˆ uj  n−1



  uj Δuj + κ ˆ uj dx



Bj

 2 Δuj + κ ˆ uj 

dx − κ ˆ n0 − 1



  uj Δuj + κ ˆ uj dx



Bj

|αj |2 φj (ˆ κ; u j , uj ) = 0 .

j=1

  This shows that Φ(ˆ κ)u, u ∗ ≤ 0 for all u ∈ Vm . By Corollary A.55 we conclude that there exist eigenvalues λ (ˆ κ) ≤ −1 of κ ˆA + κ ˆ 2 B for  = 1, . . . , m. Furthermore, again by Corollary A.55 of Appendix A.6 the eigenvalues λ (κ) depend continuously on κ and λ (0) = 0 for all . Therefore, for every  ∈ {1, . . . , m} ˆ ] with λ (κ ) = −1. The corresponding eigenfunctions there exists κ ∈ (0, κ  √ κ :  = 1, . . . , m u ,  = 1, . . . , m, satisfy Φ(κ )u = 0 which implies that are interior transmission eigenvalues. Since m was arbitrary the existence of

296

Inverse Scattering Problem 

infinitely many real eigenvalues is shown.

The—very natural—question of existence of complex transmission eigenvalues in this general situation is totally open. Since even in the case of a constant refractive index in a ball both, existence and nonexistence, of complex eigenvalues can occur (see, for example Theorem 7.48) the problem is certainly hard to solve. Aspects concerning the asymptotic distribution of the eigenvalues (Weyl’s formula), location and properties of the eigenfunctions have attracted a lot of attention, also for much more general types of elliptic problems. We only mention the special issue ([37]) of Inverse Problems in 2013, the already mentioned monographs [34, 55], and refer to the references therein.

7.6.3

The Inverse Spectral Problem for the Radially Symmetric Case

After the study of the existence and discreteness of the interior transmission eigenvalues the natural question arises where these eigenvalues determine the refractive index n(x) uniquely. For spherically stratified media studied in Subsection 7.6.1, this question is the analogue of the inverse Sturm-Liouville problem of Chapter 5 and has been subject of intensive research. It started with the work by J. McLaughlin and P. Polyakov ([190]) and was picked up in, e.g., [56–58, 46, 45]. We will follow the presentations of [34, 55]. We assume as at the beginning of Subsection 7.6.1 that n ∈ C 2 [0, 1] is positive on [0, 1]. This smoothness assumption is mainly necessary to apply the Liouville transform. We have seen in Subsection 7.6.1 that the interior transmission eigenvalue eigenvalues of (7.44a), (7.44b) are—for radially symmetric eigenfunctions—just the zeros of the characteristic function d, given by (7.87); that is, sin k − yk (1) cos k , k ∈ C , (7.102) d(k) = yk (1) k where yk solves the initial value problem yk (r) + k 2 n(r)yk (r) = 0 in (0, 1) ,

yk (0) = 0 , yk (0) = 1 .

(7.103)

Therefore, the inverse problem is to recover the refractive index n = n(r) from the zeros of the characteristic function d = d(k). This corresponds exactly to the situation of Chapter 5 where the inverse spectral problem was to recover q = q(x) from the zeros of the characteristic function f = f (λ). Therefore, it is not surprising that we use similar arguments. We saw already (in Lemma 7.44) that the characteristic function d is an 1 even entire function of order one provided n(1) = 1 or η := 0 n(s)ds = 1. We prove a further property of d. Lemma 7.53 Let d be the characteristic function from (7.102). Then d(k) 1 − = lim k→0 k 2 3

1 0

n(s) s2 ds .

7.6

The Interior Transmission Eigenvalue Problem

297

which is not zero if, for example, n(r) ≤ 1 for all r ∈ [0, 1] and n ≡ 1. Therefore, under this condition k = 0 is a root of d of order 2. Proof: For low values of |k| it is convenient to use an integral equation argument directly for yk . Indeed, it is not difficult to show that yk ∈ C 2 [0, 1] solves (7.103) if, and only if yk ∈ C[0, 1] satisfies the Volterra equation r 2 (r − s) n(s) yk (s) ds , 0 ≤ r ≤ 1 . (7.104) yk (r) = r − k 0

For sufficiently small |k| (actually, for all values of k) this fixed point equation is solved by the Neumann series (see Theorem A.31 of Appendix A). The first two terms yield r 2 (r − s) n(s) s ds + O(|k|4 ) . yk (r) = r − k 0

Also, for the derivative we get yk (r)

= 1 − k

2

r n(s) yk (s) ds = 1 − k 0

2

r

n(s) s ds + O(|k|4 ) .

0

Substituting these expansions for r = 1 and the power series of sin k/k and cos k into the definition of d yields after collecting the terms with k 2 ⎡ ⎤ 1 1 d(k) = k 2 ⎣ − n(s) s2 ds⎦ + O(|k|4 ) . 3 0

This ends the proof.



From the properties derived in Lemmas 7.44 and  1 7.53 we have the following form of the characteristic function, provided 0 n(s)s2 ds = 1/3 holds and n(1) = 1 or η = 1.  ∞

% k2 1 − 2 , k ∈ C, (7.105) d(k) = γ k 2 kj j=1 for some γ ∈ C where kj ∈ C are all nonzero roots of d repeated according to multiplicity. This follows directly from Theorem 7.45 and the fact that with k also −k is a zero. Indeed, let {kj : j ∈ J} be the set of zeros in the half plane {z ∈ C : Re z > 0 or z = it, t > 0}. Then the disjoint union {kj : j ∈ J} ∪ {−kj : j ∈ J} cover all of the roots. We group the factors of Hadamard’s formula (7.5) into pairs     2 1 − kkj ek/kj 1 + kkj e−k/kj = 1 − kk2 for j ∈ J. Furthermore, the polynomial j

p(k) must be constant because d is even. This shows that   ∞

%

k2 1 p 2% k2 e k 1− 2 = 1− 2 . d(k) = ep k 2 kj 2 kj j=1 j∈J

298

Inverse Scattering Problem

After these preparations we turn to the inverse problem to determine n(r) from (all of) the roots kj ∈ C of the characteristic function d. Our main concern is again the question of uniqueness, namely, do the refractive indices n1 and n2 have to coincide if the roots of the corresponding characteristic functions coincide? Therefore, we assume that we have two positive refractive indices n1 , n2 ∈ C 2 [0, 1] such that the zeros kj of their characteristic functions d1 and d2 coincide. 1 Furthermore, we assume that 0 n (s) s2 ds = 1/3 for  = 1, 2 and also n (1) = 1 for  = 1, 2. Then the characteristic functions d (k) have the representations (7.105) with constants γ for  = 1, 2. Then we conclude that d1 (k)/γ1 = d2 (k)/γ2 for all k ∈ C; that is, we can determine d(k)/γ from the data {kj : j ∈ N}. In the following we proceed in several steps. StepA:We determine η from the data; that is, we show that η1 = η2 where 1 η = 0 n (s) ds for  = 1, 2 if the transmission eigenvalues corresponding to n1 and n2 coincide. Fix a ∈ R and set k = a + it for t > 0. From the asymptotic behavior (7.90) we conclude that k

d(k) γ

     1  1/4 A sin k(1 + η) + B sin k(1 − η) γ n(0)n(1)

 exp(t(η + 1)) +O t   A −ia(η+1) t(η+1) 1 + O(1/t) , t → ∞ , e = −  1/4 e γ n(0)n(1) 2i

=

because 1 + η > |1 − η|. Here, A and B are given by (7.91) and A = 0 because of n(1) = 1. Analogously, we have for the complex conjugate k = a − it k

  A d(k) =  eia(η+1) et(η+1) 1 + O(1/t) , t → ∞ .  1/4 γ γ n(0)n(1) 2i

Therefore, also the ratio is known and also ψ(a) := lim

t→∞

k d(k) = −e2ia(1+η) k d(k

for all a ∈ R .

This determines η through 2i(1 + η) = ψ  (0). Step B: Under the assumption η = 1 we determine n(1) and γ n(0)1/4 from the data. Let now k > 0 be real valued. Estimate (7.90) has the form k

     1 d(k) A sin k(1 + η) + B sin k(1 − η) + R(k) =   1/4 γ γ n(0)n(1)

7.6

The Interior Transmission Eigenvalue Problem

299

with |R(k)| ≤ c1 /k for k ≥ 1 and some c1 > 0. The left hand side is known and also η from Step A. Therefore, for fixed a > 0, also

ψ1 (T )

=

1 T

T k

  d(k) sin k(1 + η) dk γ

a

=

1

 1/4 γ n(0)n(1)

1 T

T



  A sin2 k(1 + η) +

a

    1 +B sin k(1 − η) sin k(1 + η) dk + T

T

  R(k) sin k(1 + η) dk

a

and ψ2 (T )

=

1 T

T k

  d(k) sin k(1 − η) dk γ

a

=

1

 1/4 γ n(0)n(1)

1 T

T



    A sin k(1 + η) sin k(1 − η) +

a

 1 +B sin k(1 − η) dk + T 2



T

  R(k) sin k(1 − η) dk

a

are known for T ≥ a. The terms involving R(k) tend to zero as T tends T T c1 ln(T /a) . The other elemento infinity because T1 a |R(k)| dk ≤ cT1 a dk k = T  T tary integrals can be computed explicitly which yields that lim T1 a sin k(1 + T →∞      T η) sin k(1 − η) dk = 0 and lim T1 a sin2 k(1 ± η) dk = 12 (note that η = 1). T →∞

Therefore, the limits lim ψ1 (T ) = T →∞

are known and thus also

A 2γ[n(0)n(1)]1/4

B and lim ψ2 (T )= 2γ[n(0)n(1)] 1/4 T →∞

 n(1) − 1 A ψ1 (T ) = =  . T →∞ ψ2 (T ) B n(1) + 1 lim

This determines n(1), thus also A and B and therefore also γ n(0)1/4 . From now on we assume that n(r) ≤ 1 for all r ∈ [0, 1] and n(1) < 1. Then 1 η < 1 and 0 n(s) s2 ds < 13 . Therefore, all of the assumptions are satisfied for the determination of η, n(1), and γ n(0)1/4 . −1/4 zk (s(r)) from Step C: Now we use the Liouville transform  r  yk (r) = n(s(r)) (7.88a)–(7.88d) again where s(r) = 0 n(s)ds and zk (s) solves (7.88b) in (0, η) with zk (0) = 0 and zk (0) = n(0)−1/4 . The application of Theorem 5.19 in Example 5.20 yields an explicit form of n(0)1/4 zk (r) in terms of the solution

300

Inverse Scattering Problem

K ∈ C 2 (Δ0 ) of the Gousat problem11 Kss (s, t) − Ktt (s, t) − q(s) K(s, t) = 0 in Δ0 , K(s, 0) K(s, s)

(7.106a)

= 0, 0 ≤ s ≤ η, s 1 q(σ) dσ , 0 ≤ s ≤ η , = 2

(7.106b) (7.106c)

0

where q(s) is related to n(r) by (7.88c) and Δ0 = {(s, t) ∈ R2 : 0 < t < s < η}. Indeed, equation (5.52) yields ⎡ ⎤ s(r)      sin(kt) ⎥ 1 ⎢ sin ks(r) + dt⎦ (7.107) yk (r) = K s(r), t ⎣ 1/4 k k [n(0)n(r)] 0

for 0 ≤ r ≤ 1 and thus

⎡

sin(kη) 1 + yk (1) =  1/4 ⎣ k n(0)n(1)

⎤

η K(η, t) 0

sin(kt) ⎦ dt . k

Differentiation of (7.107) and setting r = 1 yields ⎤ ⎡  1/4 η η n(1) ∂K(η, t) sin(kt) ⎦ ⎣cos(kη) + sin(kη) q(s) ds + yk (1) = dt n(0) 2k ∂s k 0 0 ⎤ ⎡ η   n (1) sin(kt) ⎦ ⎣ sin(kη) + − dt K(η, t) k k 4 n(0)1/4 n(1)3/4 0

Step D: We determine K(η, t) for t ∈ [0, η]. With the constant γ from Hadamard’s formula (7.105) we have for k = π,  ∈ N, π

d(π) γ

= π =

yπ (1) (−1)+1 γ

⎡ ⎤ η (−1)+1 ⎣sin(πη) + K(η, t) sin(πt) dt⎦ . [γ n(0)1/4 ] n(1)1/4 0

The left hand  ηside and the factor in front of the bracket are known which implies that 0 K(η, t) sin(πt) dt is determined from the data for all  ∈ N. This determines K(η, ·) because {sin(πt) :  ∈ N} is complete in L2 (0, η) since η < 1. Indeed, extend K(η, ·) by zero into (0, 1) and then to an odd function 11 The

interval [0, 1] is now replaced by [0, η] which does not affect the result.

7.6

The Interior Transmission Eigenvalue Problem

301

η into (−1, 1). Then 0 K(η, t) sin(πt) dt are the Fourier coefficients of this odd extension which determine K(η, ·) uniquely. Step E: Determination of n (1). From the previous arguments yk (1)/γ is known  yk (1) yk (1) and thus also ψ(k) := k d(k) γ + k γ cos k = γ sin k. We determine the asymptotic behavior of this  η expression as k ∈ R tends to infinity. The integrals η K(η, t) sin(kt) dt and ∂K(η, t)/∂s sin(kt) dt tend to zero as O(1/k) as seen 0 0 from partial integration. Therefore, ψ(k)

d(k) yk (1) y  (1) +k cos k = k sin k γ γ γ ⎡ ⎤ η n(1)1/4 ⎣ sin(kη) cos(kη) + q(s) ds⎦ sin k 2k γ n(0)1/4

= k =

0

n (1) sin(kη) − sin k + O(1/k 2 ) . 1/4 3/4 k 4 γ n(0) n(1) The left hand side is known and also the first term on the right hand side η because 0 q(s) ds = 2K(η, η). This determines n (1) from the data. Step F: Determination of ∂K(η, t)/∂s for t ∈ [0, η]. We compute ψ  (π) from the previously defined function ψ as   (1) n(1)1/4 yπ sin(πη)  cos(π) = K(η, η) (−1) cos(πη) + ψ  (π) = 1/4 γ π γ n(0)  η ∂K(η, t) sin(πt) dt + ∂s π 0 ⎤ ⎡ η n (1) sin(πη) sin(πt) − + dt⎦ . (−1) ⎣ K(η, t) π π 4 γ n(0)1/4 n(1)3/4 0

η

From this we conclude that also 0 ∂K(η, t)/∂s sin(πt) dt is known for all  ∈ N and thus also ∂K(η, t)/∂s by the same arguments as in Step D. Step G: Determination of q = q(s) for 0 ≤ s ≤ η. We recall from Steps D and F, that K(η, t) and ∂K(η, t)/∂s are determined from the data for all t ∈ [0, η]. More precisely, if K denote the solutions of (7.106a)–(7.106c) for q = q ,  = 1, 2, then η1 = η2 =: η and K1 (η, ·) = K2 (η, ·) and ∂K1 (η, ·)/∂s = ∂K2 (η, ·)/∂s. The difference K := K1 − K2 satisfies Kss (s, t) − Ktt (s, t) − q1 (s) K(s, t) = [q1 (s) − q2 (s)] K2 (s, t) and K(·, 0) = 0 on [0, η] and K(s, s) =

1 2

s 0

in Δ0 ,

[q1 (σ) − q2 (σ)] dσ for 0 ≤ s ≤ η.

Furthermore, K(η, ·) = ∂K(η, ·)/∂s = 0 on [0, η]. Now we apply Theorem 5.18 from Chapter 5 with q = q1 , F = K2 , and f = g = 0. We observe that the pair

302

Inverse Scattering Problem

(K, q1 − q2 ) solves the homogeneous system (5.43a)–(5.43d). The uniqueness result of this theorem yields that K = 0 and q1 = q2 . This proves Part G. Step H: In this final part we have to determine n from q where their relationship is given by (7.88c). Let again n1 and n2 be two indices with the same transmission eigenvalues. Define u by   1/4 n r (s) ,

s ∈ [0, η ] = [0, η] , r where again r = r (s) is the inverse of s = s (r) = 0 n (σ) dσ for  = 1, 2. An elementary computation (using the chain rule and the derivative of the inverse function and (7.88c)) yields that u satisfies the ordinary linear differential equation u (s) =

u (s) = q (s) u (s) ,

0 ≤ s ≤ η, n (1)

with end conditions u (η) = n (1)1/4 and u (η) = 4 n(1)5/4 . From q1 = q2 and n1 (1) = n2 (1) and n1 (1) = n2 (1) and the uniqueness of this initial value problem we conclude that u1 (s) = u2 (s) for all s; that is,           s1 r1 (s) = n1 r1 (s) = n2 r2 (s) = s2 r2 (s) .     On the other hand, differentiating s r (s) = s yields s r (s) r (s) = 1 which implies that r1 = r2 , thus r1 = r2 and, finally, n1 = n2 . We summarize the result in the following theorem. Theorem 7.54 Let nj ∈ C 2 [0, 1], j = 1, 2, be positive with nj (r) ≤ 1 on [0, 1] and nj (1) < 1 such that all of the corresponding transmission eigenvalues with radially symmetric eigenfunctions coincide. Then n1 and n2 have to coincide.

7.7

Numerical Methods

In this section, we describe three types of numerical algorithms for the approximate solution of the inverse scattering problem for the determination of n and not only of the support D of n − 1. We assume—unless stated otherwise—that n ∈ L∞ (R3 ) with n(x) = 1 outside some ball B = B(0, a) of radius a > 0. The numerical methods we describe now are all based on the Lippmann– Schwinger integral equation. We define the volume potential V φ with density φ by  eik|x−y| φ(y) dy , x ∈ B . (7.108) (V φ)(x) := 4π|x − y| |y| 0 such that, if m ∈ C(Q) with m ∞ ≤ ˆ is the corresponding total field with exact far field pattern ε and u = u(x, θ) ∞ ˆ ˆ then the sequence (m , u ) constructed by the regularized x, θ) = u∞ (ˆ x, θ), u (ˆ   ˜ u ˜) ∈ XN × C Q × S 2 that algorithm (Ar ), (Br ), (Cr ) converges to some (m, satisfies the scattering problem with refraction contrast m. ˜ Its far field pattern  u ˜∞ coincides with u∞ at the points x ˆj , θˆj ∈ S 2 × S 2 , j ∈ ZN . If,  in addition, the exact solution m satisfies m ∈ XN , then the sequence m , u converges to (m, u). Proof:

We define the operator       L : C Q × S 2 × C S 2 × S 2 −→ XN × C Q × S 2

by



 L1 ρ , w − k 2 V (ui L1 ρ) .   Then L is a left inverse of T  (0, ui ) on XN × C Q × S 2 ; that is,   L T  (0, ui )(μ, v) = (μ, v) for all (μ, v) ∈ XN × C Q × S 2 .   Indeed, let (μ, v) ∈ XN × C Q × S 2 and set (w, ρ) = T  (0, ui )(μ, v), i.e., w = v + k 2 V (μui ) and ρ = W (μui ). The latter equation implies that   4π  μ(y) e−i j·y dy = − 2 ρ x ˆj , θˆj , j ∈ ZN . k L(w, ρ) :=

Q

Because μ ∈ XN , this yields μ = L1 ρ and thus L(w, ρ) = (μ, v). With the abbreviations z = m , u and R = (ui , u∞ ), we can write the regularized algorithm in the form   z+1 = z − L T (z ) − R ,  = 0, 1, 2, . . .   in the space XN × C Q × S 2 . We can now apply a general result about local convergence of the simplified Newton method (see Appendix A, Theorem A.65).   2 of This yields the existence of a unique solution ( m, ˜ u ˜ ) ∈ X × C Q × S N     , u ) to ( m, ˜ u ˜ ). L T m, ˜ u ˜ − R = 0 and linear convergence of the sequence (m     2 i ˜ u = u is equivalent to the scattering problem by The equation u ˜ + k V m˜     ˆj , θˆj = ˜ u = L1 f is equivalent to u ˜∞ x Theorem 7.12. The equation L1 W m˜   ˆj , θˆj forall j ∈ ZN . Finally, if m ∈ XN , then (m, u) satisfies L T (m, u) = u∞ x L R and thus m, ˜ u ˜ = (m, u). This proves the assertion.  We have formulated the algorithm with respect to the Lippmann–Schwinger integral equation because our analysis of existence and continuous dependence is based on this setting. There is an alternative way to formulate the simplified Newton method in terms of the original scattering problems; see [113]. We note also that our analysis can easily be modified to treat the case where only n ∈ L∞ (B). For numerical examples, we refer to [113].

7.7

Numerical Methods

7.7.2

307

A Modified Gradient Method

The idea of the numerical method proposed and numerically tested by Kleinman and van den Berg (see [164]) is to solve (7.112a), (7.112b) by a gradient-type method. For simplicity, we describe the method again in the function space setting and refer for discretization aspects to the original literature [164]. Again let B = B(0, a) contain the support of m = 1 − n. (A) Choose m0 ∈ L∞ (B), u0 ∈ L2 (B × S 2 ), and set  = 0. (B) Choose directions e ∈ L2 (B × S 2 ) and d ∈ L∞ (B), and set u+1 = u + α e ,

m+1 = m + β d .

(7.122)

The stepsizes α , β > 0 are chosen in such a way that they minimize the functional Ψ (α, β) :=

r+1 2L2 (B×S 2 )

s+1 2L2 (S 2 ×S 2 )

,

(7.123a)

where the defects r+1 and s+1 are defined by   r+1 := ui − u+1 − k 2 V m+1 u+1 ,   s+1 := u∞ − W m+1 u+1 .

(7.123b)

ui 2L2 (B×S 2 )

+

f 2L2 (S 2 ×S 2 )

(7.123c)

(C) Replace  by  + 1 and continue with step (B). There are different choices for the directions d and e . In [164],  ˆ u (x, θ) ˆ ds(θ) ˆ , x ∈ B , and e := r d (x) = − d˜ (x, θ)

(7.124)

S2

have been chosen where     d˜ = −W ∗ W (m u ) − u∞ ∈ L∞ B × S 2 . In this case, d is the steepest descent direction of μ → W (μu )−u∞ 2L2 (S 2 ×S 2 ) . In [266], for d and e Polak–Ribi`ere conjugate gradient directions are chosen. A rigorous convergence analysis of either method has not been carried out. A severe drawback of the methods discussed in Sections 7.7.1 and 7.7.2 is ˆ To estimate the that they iterate on functions m = m (x) and u = u (x, θ). storage requirements, we choose a grid of order N ·N ·N grid points in B and M directions θ1 , . . . , θM ∈ S 2 . Then both methods iterate on vectors of dimension N 6 · M . From the uniqueness results, M is expected to be large, say, of order N 2 . For large values of M , the method described next has proven to be more efficient.

308

7.7.3

Inverse Scattering Problem

The Dual Space Method

The method described here is due to Colton and Monk [61, 62] based on their earlier work for inverse obstacle scattering problems (see [59, 60]). There exist various modifications of this method, but we restrict ourselves to the simplest case. This method consists of two steps. In the first step, one tries to determine ˆ such that the corresponding a superposition of the incident fields ui = ui (·, θ) ˆ far field pattern u∞ (·, θ) is (close to) the far field pattern of radiating multipoles. In the second step, the function m = n − 1 is determined from an interior transmission problem. We describe both steps separately. Assume for the following that the origin is contained in B = B(0, a). By u∞ we denote again the measured far field pattern. Step 1: Determine g ∈ L2 (S 2 ) with  ˆ g(θ) ˆ ds(θ) ˆ = 1, u∞ (ˆ x, θ)

x ˆ ∈ S2 .

(7.125)

S2

ˆ x, θ) In Theorem 7.22, we have proven that for the exact far field pattern u∞ (ˆ 2 2 this integral equation of the first kind is solvable in L (S ) if and only if the interior transmission problem Δv + k 2 v = 0 in B ,

Δw + k 2 nw = 0 in B ,

w(x) − v(x) = ∂w(x) ∂v(x) − ∂ν ∂ν

=

(7.126a)

ik|x|

e

|x|

on ∂B ,

∂ eik|x| ∂ν |x|

on ∂B ,

(7.126b) (7.126c)

has a solution v, w ∈ L2 (B) in the ultra weak sense of Definition 7.21 such that  v(x) = eikx·ˆy g(ˆ y ) ds(ˆ y ) , x ∈ R3 . (7.127) S2

The kernel of the integral operator in (7.125) is (for the exact far field pattern) analytic with respect to both variables, thus (7.125) represents a severely illposed—but linear—equation and can be treated by Tikhonov’s regularization method as described in Chapter 2 in detail. (In this connection, see the remark following Theorem 7.23.) We formulate the interior transmission problem (7.126a)–(7.126c) as an integral equation. Lemma 7.56 (a)Let v, w ∈ L2 (B) solve the boundary value problem (7.126a)– (7.126c) in the ultra weak sense of (7.46b). Define u = w − v in B and

7.7

Numerical Methods

309

2 u(x) = exp(ik|x|)/|x| in R3 \ B. Then u ∈ Hloc (R3 ) and u and w ∈ L2 (B) solve  2 m(y) w(y) Φ(x, y) dy , x ∈ R3 , (7.128a) u(x) = k |y| 0 and λ1 , λ2 > 0, minimize J(g, w, m) on L2 (S 2 ) × L2 (B) × C,

(7.130a)

7.7

Numerical Methods

311

where J(g, w, m)

:= F g − 1 2L2 (S 2 ) + ε g 2L2 (S 2 ) 2

+ λ1 w − vg − k V

(7.130b)

(mw) 2L2 (B)

+ λ2 k 2 V (mw) − 4π Φ(·, 0) 2L2 (∂B) , and the far field operator F : L2 (S 2 ) → L2 (S 2 ) is defined by (see (7.40))        ˆ , x ˆ, θˆ g θˆ ds(θ) ˆ ∈ S2. u∞ x (F g) x ˆ := S2

Theorem 7.57 This optimization problem (7.130a), (7.130b) has an optimal solution (g, w, m) for every choice of ε, λ1 , λ2 > 0 and every compact subset C ⊂ L∞ (B). Proof: Let (gj , wj , mj ) ∈ L2 (S 2 ) × L2 (B) × C be a minimizing sequence; that is, J(gj , wj , mj ) → J ∗ where the optimal value J ∗ is defined by   J ∗ := inf J(g, w, m) : (g, w, m) ∈ L2 (S 2 ) × L2 (B) × C . We can assume that (mj ) converges to some m ∈ C because C is compact. Several tedious applications of the parallelogram equality

a + b 2 = − a − b 2 + 2 a 2 + 2 b 2 and the binomial formula

b 2 = a 2 + 2 Re (a, b − a) + a − b 2 yield −J

∗



 1 1 (gj + g ), (wj + w ), mj ≥ −J 2 2 1 1 = − J(gj , wj , mj ) − J(g , w , mj ) 2 2 1 ε + F (gj − g ) 2L2 (S 2 ) + gj − g 2L2 (S 2 ) 4 4   λ1

(wj − w ) − vgj −g − k 2 V mj (wj − w ) 2L2 (B) + 4   λ2 k 4

V mj (wj − w ) 2L2 (∂B) . + 4

From this we conclude that 1 1 J(gj , wj , mj ) + J(g , w , mj ) 2 2   λ1 ε

gj − g 2L2 (S 2 ) +

(wj − w ) − vgj −g − k 2 V mj (wj − w ) 2L2 (B) . 4 4

−J ∗ + ≥

312

Inverse Scattering Problem

The left-hand side tends to zero as j and  tend to infinity, therefore we conclude that (gj ) is a Cauchy sequence, thus converging gj → g in L2 (S 2 ). Furthermore, from  

(wj − w ) − vgj −g − k 2 V mj (wj − w ) L2 (B) −→ 0 as , j → ∞ and the convergence gj → g we conclude that  

(wj − w ) − k 2 V mj (wj − w ) L2 (B) −→ 0 as , j → ∞. The operators I − k 2 V (mj ·) converge to the isomorphism I − k 2 V (m·) in the operator norm of L2 (B). Therefore, by Theorem A.37 of Appendix A, we conclude that (wj ) is a Cauchy sequence and thus is convergent in L2 (B) to some w. The continuity of J implies that J(gj , wj , mj ) → J(g, w, m). Therefore, (g, w, m) is optimal. 

7.8

Problems

p 7.1 Let Q = (−π, π)3 ⊂ R3 be the cube and H p (Q), H0p (Q), and Hper (Q) be the Sobolev spaces defined at the beginning of this chapter. Show that p p (Q) ⊂ H p (Q) and H0p (Q) ⊂ Hper (Q) with bounded inclusions. Use Hper p this result to show that H0 (Q) is compactly imbedded in L2 (Q).

ˆ k) and ub (ˆ ˆ x, θ, 7.2 Let ub1,∞ (ˆ 2,∞ x, θ, k) be the far field patterns of the Born ˆ wave approximations corresponding to observation x ˆ, angle of incidence θ, number k, and indices of refraction n1 and n2 , respectively. Assume that ˆ k) = ub (ˆ ˆ ub1,∞ (ˆ x, θ, 2,∞ x, θ, k) for all x ˆ ∈ S 2 and k ∈ [k1 , k2 ] ⊂ R+ and some θˆ ∈ S 2 . Prove that n1 = n2 .   ˆ; θˆ , 7.3 Prove the following result, sometimes called Karp’s theorem. Let u∞ x x ˆ, θˆ ∈ S 2 , be the far field pattern and assume that there exists a function f : [−1, 1] → C with     ˆ; θˆ = f x ˆ · θˆ for all x ˆ, θˆ ∈ S 2 . u∞ x Prove that the index of refraction n has to be radially symmetric: n = n(r). Hint: Rotate the geometry and use the uniqueness result. 7.4 Show that for any a > 0 max



|x|≤a |y| a, and solve this elliptic equation explicitly by separation of variables.

7.8

Problems

313

7.5 Show that the characteristic functions fη from (7.94) for η = and η = 2/3 have the forms f1/2 (k)

=

f2/3 (k)

=

k , k ∈ C, 3   2 2k 3 k sin 3 + 2 cos , 3 3 3

√ n = 1/2

sin3

k ∈ C,

respectively. Discuss the existence of zeros of these functions and justify for these examples the assertions of Theorem 7.48. Hint: Use the addition formulas for the trigonometric functions to express fη in terms of sin(k/3) and cos(k/3).

Appendix A

Basic Facts from Functional Analysis In this appendix, we collect some of the basic definitions and theorems from functional analysis. We prove only those theorems whose proofs are not easily accessible. We recommend the monographs [151, 168, 230, 271] for a comprehensive treatment of linear and nonlinear functional analysis.

A.1

Normed Spaces and Hilbert Spaces

First, we recall two basic definitions. Definition A.1 (Scalar Product, Pre-Hilbert Space) Let X be a vector space over the field K = R or K = C. A scalar product or inner product is a mapping (·, ·)X : X × X −→ K with the following properties: (i) (x + y, z)X = (x, z)X + (y, z)X for all x, y, z ∈ X, (ii) (αx, y)X = α (x, y)X for all x, y ∈ X and α ∈ K, (iii) (x, y)X = (y, x)X for all x, y ∈ X, (iv) (x, x)X ∈ R and (x, x)X ≥ 0, for all x ∈ X, (v) (x, x)X > 0 if x = 0. A vector space X over K with inner product (·, ·)X is called a pre-Hilbert space over K. © Springer Nature Switzerland AG 2021 A. Kirsch, An Introduction to the Mathematical Theory of Inverse Problems, Applied Mathematical Sciences 120, https://doi.org/10.1007/978-3-030-63343-1

315

316

Appendix A The following properties are easily derived from the definition:

(vi) (x, y + z)X = (x, y)X + (x, z)X for all x, y, z ∈ X, (vii) (x, αy)X = α (x, y)X for all x, y ∈ X and α ∈ K. Definition A.2 (Norm) Let X be a vector space over the field K = R or K = C. A norm on X is a mapping  · X : X −→ R with the following properties: (i) xX > 0 for all x ∈ X with x = 0, (ii) αxX = |α| xX for all x ∈ X and α ∈ K, (iii) x + yX ≤ xX + yX for all x, y ∈ X. A vector space X over K with norm  · X is called normed space over K. Property (iii) is called triangle inequality. Applying it to the identities x = (x − y) + y and y = (y − x) + x yields the second triangle inequality x − yX ≥  xX − yX  for all x, y ∈ X. Theorem A.3 Let X be a pre-Hilbert space. The mapping  · X : X −→ R defined by  (x, x)X , x ∈ X , xX := is a norm; that is, it has properties (i), (ii), and (iii) of Definition A.2. Furthermore, (iv) |(x, y)X | ≤ xX yX for all x, y ∈ X (Cauchy–Schwarz inequality), (v) x ± y2X = x2X + y2X ± 2 Re(x, y)X for all x, y ∈ X (binomial formula), (vi) x + y2X + x − y2X = 2x2X + 2y2X for all x, y ∈ X. In the following example, we list some of the most important pre-Hilbert and normed spaces. Example A.4 (a) Cn is a pre-Hilbert space of dimension n over C with inner product n (x, y)2 := k=1 xk y k . (b) Cn is a pre-Hilbert space of dimension 2n over R with inner product n (x, y)2 := Re k=1 xk y k . (c) Rn is a pre-Hilbert space of dimension n over R with inner product n (x, y)2 := k=1 xk yk .

A.1 Normed Spaces and Hilbert Spaces

317

(d) For p ≥ 1 define the set p of complex-valued sequences by   ∞  p p |xk | < ∞ .  := (xk ) :

(A.1)

k=1

Then p is a linear space over C because if (xk ), (yk ) ∈ p , then (λxk ) and (xk + yk ) are also in p . The latter follows from the inequality |xk + yk |p ≤ (2 max{|xk |, |yk |})p ≤ 2p (|xk |p + |yk |p ).  xp :=

∞ 

1/p |xk |

p

,

x = (xk ) ∈ p ,

k=1

defines a norm in p . The triangle inequality in the case p > 1 is known as the Minkowski inequality. In the case p = 2, the sesquilinear form (x, y)2 :=

∞ 

x = (xk ) , y = (yk ) ∈ 2 ,

xk yk ,

k=1

defines an inner product on 2 . It is well-defined by the Cauchy–Schwarz inequality. (e) The space C[a, b] of (real- or complex-valued) continuous functions on [a, b] is a pre-Hilbert space over R or C with inner product

b (x, y)L2 :=

x(t) y(t) dt ,

x, y ∈ C[a, b] .

(A.2a)

a

The corresponding norm is called the Euclidean norm and is denoted by 

b  (x, x)L2 = |x(t)|2 dt , x ∈ C[a, b] . (A.2b) xL2 := a

(f) On the same vector space C[a, b] as in example (e), we introduce a norm by (A.3) x∞ := max |x(t)| , x ∈ C[a, b] , a≤t≤b

that we call the supremum norm. (g) Let m ∈ N and α ∈ (0, 1]. We define the spaces C m [a, b] and C m,α [a, b] by   x is m times continuously m C [a, b] := x ∈ C[a, b] : , differentiable on [a, b]

C

m,α

[a, b]

  |x(m) (t) − x(m) (s)| m := x ∈ C [a, b] : sup 0 with M ⊂ B(0, r). The set M ⊂ X is called open if for every x ∈ M there exists ε > 0 such that B(x, ε) ⊂ M . The set M ⊂ X is called closed if the complement X \ M is open. (b) A sequence (xk )k in X is called bounded if there exists c > 0 such that xk X ≤ c for all k. The sequence (xk )k in X is called convergent if there exists x ∈ X such that x − xk X converges to zero in R. We denote the limit by x = limk→∞ xk , or we write xk → x as k → ∞. The sequence (xk )k in X is called a Cauchy sequence if for every  > 0 there exists N ∈ N with xm − xk X <  for all m, k ≥ N . (c) Let (xk )k be a sequence in X. A point x ∈ X is called an accumulation point if there exists a subsequence (akn )n that converges to x. (d) A set M ⊂ X is called compact if every sequence in M has an accumulation point in M . Example A.6 Let X = C[0, 1] over R and xk (t) = tk , t ∈ [0, 1], k ∈ N. The sequence (xk )k converges to zero with respect to the Euclidean norm ·L2 introduced in (A.2b). With respect to the supremum norm  · ∞ of (A.3), however, the sequence does not converge to zero. It is easy to prove (see Problem A.1) that a set M is closed if and only if the limit of every convergent sequence (xk )k in M also belongs to M . The sets int(M ) := {x ∈ M : there exists ε > 0 with B(x, ε) ⊂ M } and closure(M ) :=



x ∈ X : there exists (xk )k in M with x = lim xk k→∞

A.1 Normed Spaces and Hilbert Spaces

319

are called the interior and closure, respectively, of M . The set M ⊂ X is called dense in X if closure(M ) = X. In general, the topological properties depend on the norm in X as we have seen already in Example A.6. For finite-dimensional spaces, however, these properties are independent of the norm. This is seen from the following theorem. Theorem A.7 Let X be a finite-dimensional space with norms  · 1 and  · 2 . Then both norms are equivalent; that is, there exist constants c2 ≥ c1 > 0 with c1 x1 ≤ x2 ≤ c2 x1

for all x ∈ X .

In other words, every ball with respect to  · 1 contains a ball with respect to  · 2 and vice versa. Further properties are collected in the following theorem. Theorem A.8 Let X be a normed space over K and M ⊂ X be a subset. (a) M is closed if and only if M = closure(M ), and M is open if and only if M = int(M ). (b) If M = X is a linear subspace, then int(M ) = ∅, and closure(M ) is also a linear subspace. (c) In finite-dimensional spaces, every subspace is closed. (d) Every compact set is closed and bounded. In finite-dimensional spaces, the reverse is also true (Theorem of Bolzano–Weierstrass): In a finitedimensional normed space, every closed and bounded set is compact. A crucial property of the set of real numbers is its completeness. It is also a necessary assumption for many results in functional analysis. Definition A.9 (Banach Space, Hilbert Space) A normed space X over K is called complete or a Banach space if every Cauchy sequence converges in X. A complete pre-Hilbert space is called a Hilbert space. The spaces Cn and Rn are Hilbert spaces with respect to their canonical inner products. The space C[a, b] is not complete with respect to the inner product (·, ·)L2 of (A.2a)! As an example, we consider the sequence xk (t) = tk for 0 ≤ t ≤ 1 and xk (t) = 1 for 1 ≤ t ≤ 2. Then (xk )k is a Cauchy sequence in C[0, 2] but does not converge in C[0, 2] with respect to (·, ·)L2 because it converges to the function  0, t < 1, x(t) = 1, t ≥ 1,   that is not continuous. The space C[a, b],  · ∞ , however, is a Banach space. Every normed space or pre-Hilbert space X can be “completed”; that is, ˜ respectively, that extends there exists a “smallest” Banach or Hilbert space X, X (that is, xX = xX˜ or (x, y)X = (x, y)X˜ , respectively, for all x, y ∈ X). More precisely, we have the following (formulated only for normed spaces).

320

Appendix A

Theorem A.10  Let X be a normed space with norm  · X . There exist a ˜ such that ˜  ·  ˜ and an injective linear operator J : X → X Banach space X, X ˜ is dense in X, ˜ and (i) The range J(X) ⊂ X (ii) JxX˜ = xX for all x ∈ X; that is, J preserves the norm. ˜ is uniquely determined in the sense that if X ˆ is a second space Furthermore, X ˆ then with properties (i) and (ii) with respect to a linear injective operator J, ˆ the operator Jˆ J −1 : J(X) → J(X) has an extension to a norm-preserving ˜ onto X. ˆ In other words, X ˜ and X ˆ can be identified. isomorphism from X   We denote the completion of the pre-Hilbert space C[a, b], (·, ·)L2 by L2 (a, b). Using Lebesgue integration theory, it can be shown that the space L2 (a, b) is characterized as follows. (The notions “measurable,” “almost everywhere” (a.e.), and “integrable” are understood with respect to the Lebesgue measure.) First, we define the vector space

L2 (a, b) := x : (a, b) → C : x is measurable and |x|2 integrable , where addition and scalar multiplication are defined pointwise almost everywhere. Then L2 (a, b) is a vector space because, for x, y ∈ L2 (a, b) and α ∈ C, x + y and αx are also measurable and αx, x + y ∈ L2 (a, b), the latter by the binomial theorem |x(t) + y(t)|2 ≤ 2|x(t)|2 + 2|y(t)|2 . We define a sesquilinear form on L2 (a, b) by

b x, y :=

x(t) y(t) dt ,

x, y ∈ L2 (a, b) .

a

·, · is not an inner product on L2 (a, b) because x, x = 0 only implies that x vanishes almost everywhere; that is, that x ∈ N , where N is defined by

N := x ∈ L2 (a, b) : x(t) = 0 a.e. on (a, b) . Now we define L2 (a, b) as the factor space L2 (a, b) := L2 (a, b)/N and equip L2 (a, b) with the inner product 

 [x], [y] L2 :=

b x(t) y(t) dt ,

x ∈ [x] , y ∈ [y] .

a 2

Here, [x], [y] ∈ L (a, b) are equivalence classes of functions in L2 (a, b). Then it can be shown that this definition is well defined and yields an inner product on L2 (a, b). From now on, we write x ∈ L2 (a, b) instead of x ∈ [x] ∈ L2 (a, b). Furthermore, it can be shown by fundamental results of Lebesgue integration theory that L2 (a, b) is complete; that is, a Hilbert space and contains C[a, b] as a dense subspace.

A.2 Orthonormal Systems

321

Definition A.11 (Separable Space) The normed space X is called separable if there exists a countable dense subset M ⊂ X; that is, if there exist M and a bijective mapping j : N → M with closure(M ) = X. The spaces Cn , Rn , L2 (a, b), and C[a, b] are all separable. For the first two examples, let M consist of all vectors with rational coefficients; for the latter examples, take polynomials with rational coefficients. Definition A.12 (Orthogonal Complement) Let X be a pre-Hilbert space (over K = R or C). (a) Two elements x and y are called orthogonal if (x, y)X = 0. (b) Let M ⊂ X be a subset. The set

M ⊥ := x ∈ X : (x, y)X = 0 for all y ∈ M is called the orthogonal complement of M . ⊥  M ⊥ is always a closed subspace and M ⊂ M ⊥ . Furthermore, A ⊂ B implies that B ⊥ ⊂ A⊥ . The following theorem is a fundamental result in Hilbert space theory and relies heavily on the completeness property. Theorem A.13 (Projection Theorem) Let X be a pre-Hilbert space and V ⊂ X be a complete subspace. Then V =  ⊥ ⊥ V . Every x ∈ X possesses a unique decomposition of the form x = v + w, where v ∈ V and w ∈ V ⊥ . The operator P : X → V , x → v, is called the orthogonal projection operator onto V and has the properties (a) P v = v for v ∈ V ; that is, P 2 = P ; (b) x − P xX ≤ x − v  X for all v  ∈ V . This means that P x ∈ V is the best approximation of x ∈ X in the subspace V.

A.2

Orthonormal Systems

In this section, let X always be a separable Hilbert space over the field K = R or C. Definition A.14 (Orthonormal System) A countable set of elements A = {xk : k = 1, 2, 3, . . .} is called an orthonormal system (ONS) if (i) (xk , xj )X = 0 for all k = j and

322

Appendix A

(ii) xk X = 1 for all k ∈ N. A is called a complete or a maximal orthonormal system if, in addition, there is no ONS B with A ⊂ B and A = B. One can show using Zorn’s Lemma that every separable Hilbert possesses a maximal ONS. Furthermore, it is well known from linear algebra that every countable set of linearly independent elements of X can be orthonormalized. For any set A ⊂ X, let  span A :=

n 

 αk xk : αk ∈ K, xk ∈ A, n ∈ N

(A.5)

k=1

be the subspace of X spanned by A. Theorem A.15 Let A = {xk : k = 1, 2, 3, . . .} be an orthonormal system. Then (a) Every finite subset of A is linearly independent. (b) If A is finite; that is, A = {xk : k = 1, 2, . . . , n}, then for every x ∈ X there exist uniquely determined coefficients αk ∈ K, k = 1, . . . , n, such that   n      αk xk  ≤ x − aX for all a ∈ span A . (A.6) x −   k=1

X

The coefficients αk are given by αk = (x, xk )X for k = 1, . . . , n. (c) For every x ∈ X, the following Bessel inequality holds: ∞    (x, xk )X 2 ≤ x2X ,

(A.7)

k=1

and the series

∞

k=1 (x, xk )X xk

converges in X.

(d) A is complete if and only if span A is dense in X. (e) A is complete if and only if for all x ∈ X the following Parseval equation holds: ∞  |(x, xk )X |2 = x2X . (A.8) k=1

(f ) A is complete if and only if every x ∈ X has a (generalized) Fourier expansion of the form x =

∞ 

(x, xk )X xk ,

k=1

(A.9)

A.2 Orthonormal Systems

323

where the convergence is understood in the norm of X. In this case, the Parseval equation holds in the following more general form: (x, y)X =

∞ 

(x, xk )X (y, xk )X .

(A.10)

k=1

This important theorem includes, as special examples, the classical Fourier expansion of periodic functions and the expansion with respect to orthogonal polynomials. We recall two examples. Example A.16 (Fourier Expansion) √ (a) The functions xk (t) := exp(ikt)/ 2π, k ∈ Z, form a complete system of orthonormal functions in L2 (0, 2π). By part (f) of the previous theorem, every function x ∈ L2 (0, 2π) has an expansion of the form

2π ∞ 1  ikt x(t) = e x(s) e−iks ds , 2π k=−∞

0

where the convergence is understood in the sense of L2 ; that is, 2

2π

2π N    1 ikt −iks  x(t) − e x(s) e ds   dt −→ 0 2π k=−M

0

0

as M, N tend to infinity. Parseval’s identity holds the form  k∈Z

|ak |

2

1 x2L2 , = 2π

ak

1 = 2π

2π

x(s) e−iks ds .

(A.11)

0

For smooth periodic functions, one can even show uniform convergence (see Section A.4). (b) The Legendre polynomials Pk , k = 0, 1, . . ., form a maximal orthonormal system in L2 (−1, 1). They are defined by Pk (t) = γk

dk (1 − t2 )k , dtk

t ∈ (−1, 1) , k ∈ N0 ,

with normalizing constants  γk =

2k + 1 1 . 2 k! 2k

We refer to [135] for details. Other important examples will be given later.

324

A.3

Appendix A

Linear Bounded and Compact Operators

For this section, let X and Y always be normed spaces and A : X → Y be a linear operator. Definition A.17 (Boundedness, Norm of A) The linear operator A is called bounded if there exists c > 0 such that AxY ≤ cxX

for all x ∈ X .

The smallest of these constants is called the norm of A; that is, AL(X,Y ) := sup x=0

AxY . xX

(A.12)

Theorem A.18 The following assertions are equivalent: (a) A is bounded. (b) A is continuous at x = 0; that is, xj → 0 implies that Axj → 0. (c) A is continuous for every x ∈ X. The space L(X, Y ) of all linear bounded mappings from X to Y with the operator norm is a normed space; that is, the operator norm has properties (i), (ii), and (iii) of Definition A.2 and the following: Let B ∈ L(X, Y ) and A ∈ L(Y, Z); then A B ∈ L(X, Z) and A BL(X,Z) ≤ AL(Y,Z) BL(X,Y ) . Integral operators are the most important examples for our purposes.   Theorem A.19 (a) Let k ∈ L2 (c, d) × (a, b) . The operator

b (Ax)(t) :=

k(t, s) x(s) ds ,

t ∈ (c, d) ,

x ∈ L2 (a, b) ,

(A.13)

a

is well-defined, linear, and bounded from L2 (a, b) into L2 (c, d). Furthermore,

AL(L2 (a,b),L2 (c,d))

  d b 

 ≤  |k(t, s)|2 ds dt . c

a

(b) Let k be continuous on [c, d] × [a, b]. Then A is also well-defined, linear, and bounded from C[a, b] into C[c, d] and

b A∞ := AL(C[a,b],C(c,d]) = max

|k(t, s)| ds .

t∈[c,d] a

A.3 Linear Bounded and Compact Operators

325

We can extend this theorem to integral operators with weakly singular kernels. We recall that a kernel k is called weakly singular on [a, b] × [a, b] if k is defined and continuous for all t, s ∈ [a, b], t = s, and there exist constants c > 0 and α ∈ [0, 1) such that |k(t, s)| ≤ c |t − s|−α

for all t, s ∈ [a, b] , t = s .

Theorem A.20 Let k be weakly singular on [a, b]. Then the integral operator A, defined by (A.13) for [c, d] = [a, b], is well-defined and bounded as an operator in L2 (a, b) as well as in C[a, b]. For the special case Y = K, we denote by X ∗ := L(X, K) the dual space of X. Often we write , x X ∗ ,X instead of (x) for  ∈ X ∗ and x ∈ X and call ·, · X ∗ ,X the dual pairing. The dual pairing is a bilinear form from X ∗ × X into K. The space X ∗∗ = (X ∗ )∗ is called the bidual of X. The canonical embedding J : X → X ∗∗ , defined by (Jx) := , x X ∗ ,X ,

x ∈ X ,  ∈ X∗ ,

(A.14)

is linear, bounded, one-to-one, and satisfies JxX ∗∗ = xX for all x ∈ X. We recall some important examples of dual spaces (where we write ·, · for the dual pairing). Example A.21 Let again K = R or K = C. with Kn itself. The identification I : (a) The dual of Kn can be identified  n Kn → (Kn )∗ is given by Ix, y = j=1 xj yj for x, y ∈ Kn . (b) Let p > 1 and q > 1 with p1 + 1q = 1. The dual of p (see Example A.4) can be identified with q . The identification I : q → (p )∗ is given by Ix, y =

∞ 

xj yj

for x = (xj ) ∈ q and y = (yj ) ∈ p .

j=1

(c) The dual (1 )∗ of the space 1 can be identified with the space ∞ of bounded sequences (equipped with the supremum norm). The identification I : ∞ → (1 )∗ is given by the form as in (b) for x = (xj ) ∈ ∞ and y = (yj ) ∈ 1 . (d) Let c0 ⊂ ∞ be the space of sequences in K which converge to zero, equipped with the supremum norm. Then c∗0 can be identified with 1 . The identification I : 1 → c∗0 is given by the form as in (b) for x ∈ 1 and y ∈ c0 . Definition A.22 (Reflexive Space) The normed space X is called reflexive if the canonical embedding of X into X ∗∗ is surjective; that is, a norm-preserving isomorphism from X onto the bidual space X ∗∗ .

326

Appendix A

The spaces p for p > 1 of Example A.21 are typical examples of reflexive spaces. The spaces 1 , ∞ , and c0 fail to be reflexive. The following important result gives a characterization of X ∗ in Hilbert spaces. Theorem A.23 (Riesz) Let X be a Hilbert space. For every x ∈ X, the functional x (y) := (y, x)X , y ∈ X, defines a linear bounded mapping from X to K; that is, x ∈ X ∗ . Furthermore, for every  ∈ X ∗ there exists one and only one x ∈ X with (y) = (y, x)X for all y ∈ X and X ∗ := sup y=0

|, y | = xX . yX

This theorem implies that every Hilbert space is reflexive. It also yields the existence of a unique adjoint operator for every linear bounded operator A : X −→ Y . We recall that for any linear and bounded operator A : X → Y between normed spaces X and Y the dual operator A∗ : Y ∗ → X ∗ is defined as T ∗  =  ◦ A for all  ∈ Y ∗ . Here  ◦ A is the composition of A and ; that is, ( ◦ A)x = (Ax) for x ∈ X. Theorem A.24 (Adjoint Operator) Let A : X −→ Y be a linear and bounded operator between Hilbert spaces. Then there exists one and only one linear bounded operator A∗ : Y −→ X with the property (Ax, y)Y = (x, A∗ y)X for all x ∈ X , y ∈ Y . This operator A∗ : Y −→ X is called the adjoint operator to A. For X = Y , the operator A is called self-adjoint if A∗ = A. Example A.25   (a) Let X = L2 (a, b), Y = L2 (c, d), and k ∈ L2 (c, d) × (a, b) . The adjoint A∗ of the integral operator

b k(t, s) x(s) ds ,

(Ax)(t) =

t ∈ (c, d) ,

x ∈ L2 (a, b) ,

t ∈ (a, b) ,

y ∈ L2 (c, d) .

a

is given by (A∗ y)(t) =

d k(s, t) y(s) ds , c

(b) Let the space X = C[a, b] of continuous function over C be supplied with the L2 -inner product. Define f, g : C[a, b] → C by

b x(t) dt

f (x) := a

and g(x) := x(a)

for x ∈ C[a, b] .

A.3 Linear Bounded and Compact Operators

327

Both f and g are linear. f is bounded but g is unbounded. There is an extension of f to a linear bounded functional (also denoted by f ) on L2 (a, b); that is, f ∈ L2 (a, b)∗ . By Theorem A.23, we can identify L2 (a, b)∗ with L2 (a, b) itself. For the given f , the representation function is just the constant function 1 because f (x) = (x, 1)L2 for x ∈ L2 (a, b). The adjoint of f is calculated by

b f (x) · y =

x(t) y dt = (x, y)L2 =

  x, f ∗ (y) L2

a 2

for all x ∈ L (a, b) and y ∈ C. Therefore, f ∗ (y) ∈ L2 (a, b) is the constant function with value y. (c) Let X be the Sobolev space H 1 (a, b); that is, the space of L2 -functions that possess generalized L2 -derivatives:   there exists α ∈ K and y ∈ L2 (a, b) with 1 2  . H (a, b) := x ∈ L (a, b) : t x(t) = α + a y(s) ds for t ∈ (a, b) We denote the generalized derivative y ∈ L2 (a, b) by x . We observe that H 1 (a, b) ⊂ C[a, b] with bounded embedding. As an inner product in H 1 (a, b), we define (x, y)H 1 := x(a) y(a) + (x , y  )L2 ,

x, y ∈ H 1 (a, b) .

Now let Y = L2 (a, b) and A : H 1 (a, b) −→ L2 (a, b) be the operator x → x for x ∈ H 1 (a, b). Then A is well-defined, linear, and bounded. It is easily seen that the adjoint of A is given by ∗

t

(A y)(t) =

y(s) ds,

t ∈ (a, b) ,

y ∈ L2 (a, b) .

a

In the following situation, we consider the case that a Banach space V is contained in a Hilbert space X with bounded imbedding j : V → X such that also j(V ) is dense in X. We have in mind examples such as H 1 (0, 1) ⊂ L2 (0, 1). Then the dual operator j ∗ is a linear bounded operator from X ∗ into V ∗ with dense range (the latter follows from the injectivity of j). Also, j ∗ is one-to-one because j(V ) is dense in X, see Problem A.3. Now we use the fact that X and X ∗ are anti-isomorphic by the Theorem A.23 of Riesz; that is, the operator jR : X  x → x ∈ X ∗ where x (z) = (z, x)X for z ∈ X is bijective and anti-linear; that is, satisfies jR (λx + μy) = λjR x + μjR y for all x, y ∈ X and λ, μ ∈ K. Therefore, also the composition j ∗ ◦ jR : X → V ∗ is anti-linear. For this reason we define the anti-dual space V  of V by

(A.15) V = V∗ = :V →K:∈V∗ where (v) = (v) for all v ∈ V . Then the operator j  := j ∗ ◦ jR : X → V  is linear and one-to-one with dense range, thus an imbedding. In this sense, X is

328

Appendix A

densely imbedded in V  . We denote the application of  ∈ V  to v ∈ V by , v and note that (, v) → , v is a sesquilinear form on V  × V . (It should not be mixed up with the dual pairing ·, · V ∗ ,V : V ∗ × V → K which is bilinear.) From this analysis, we conclude that (x, v)X = x (v) = j  x, v for all x ∈ X and v ∈ V and thus   (x, v)X  ≤ j  xV  vV for all x ∈ X , v ∈ V . (A.16) Definition A.26 (a) A Gelfand triple (or rigged Hilbert space, see [100]) V ⊂ X ⊂ V  consists of a reflexive Banach space V , a separable Hilbert space X, and the anti-dual space V  of V (all over the same field K = R or K = C) such that V is a dense subspace of X, and the imbedding j : V → X is bounded. Furthermore, the sesquilinear form ·, · : V  × V → K is an extension of the inner product in X; that is, x, v = (x, v)X for all v ∈ V and x ∈ X. (b) A linear bounded operator K : V  → V is called coercive if there exists γ > 0 with   x, Kx  ≥ γ x2V  for all x ∈ V  . (A.17) The operator K satisfies G˚ arding’s inequality if there exists a linear compact operator C : V  → V such that K + C is coercive; that is,   x, (K + C)x  ≥ γ x2V  for all x ∈ V  . By the same arguments as in the proof of the Lax–Milgram theorem (see [129]), it can be shown that every coercive operator is an isomorphism from V  onto V . Coercive operators play an important role in the study of partial differential equations and integral equations by variational methods. Often, the roles of V and V  are interchanged. For integral operators that are “smoothing”, our definition seems more appropriate. However, both definitions are equivalent in the sense that the inverse operator K −1 : V → V  is coercive in the usual sense with γ replaced by γ/K2L(V  ,V ) . The following theorems are two of the most important results of linear functional analysis. Theorem A.27 (Open Mapping Theorem) Let X, Y be Banach spaces and A : X → Y a linear bounded operator from X onto Y . Then A is open; that is, the images A(U ) ⊂ Y are open in Y for all open sets U ⊂ X. In particular, if A is a bounded isomorphism from X onto Y , then the inverse A−1 : Y → X is bounded. This result is sometimes called the Banach–Schauder theorem. Theorem A.28 (Banach–Steinhaus, Principle of Uniform Boundedness) Let X be a Banach space, Y be a normed space, I be an index set, and Aα ∈ L(X, Y ), α ∈ I, be a collection of linear bounded operators such that sup Aα xY < ∞ α∈I

Then supα∈I Aα L(X,Y ) < ∞.

for every x ∈ X .

A.3 Linear Bounded and Compact Operators

329

As an immediate consequence, we have the following. Theorem A.29 Let X be a Banach space, Y be a normed space, D ⊂ X be a dense subspace, and An ∈ L(X, Y ) for n ∈ N. Then the following two assertions are equivalent: (i) An x → 0 as n → ∞ for all x ∈ X. (ii) supn∈N An L(X,Y ) < ∞ and An x → 0 as n → ∞ for all x ∈ D. We saw in Theorem A.10 that every normed space X possesses a unique com˜ Every linear bounded operator defined on X can also be extended pletion X. ˜ to X. ˜ Y˜ be Banach spaces, X ⊂ X ˜ a dense subspace, and Theorem A.30 Let X, ˜ A : X → Y be linear and bounded. Then there exists a linear bounded operator ˜ → Y˜ with A˜ : X ˜ = Ax for all x ∈ X; that is, A˜ is an extension of A, and (i) Ax ˜ (ii) A ˜ Y˜ ) = AL(X,Y ) . L(X, Furthermore, the operator A˜ is uniquely determined. We now study equations of the form x − Kx = y ,

(A.18)

where the operator norm of the linear bounded operator K : X → X is small. The following theorem plays an essential role in the study of Volterra integral equations. Theorem A.31 (Contraction Theorem, Neumann Series) Let X be a Banach space over R or C and K : X → X be a linear bounded operator with 1/n

lim sup K n L(X) < 1 .

(A.19)

n→∞

∞ Then I −K is invertible, the Neumann series n=0 K n converges in the operator norm, and ∞  K n = (I − K)−1 . n=0

Condition (A.19) is satisfied if, for example, K m L(X) < 1 for some m ∈ N. Example A.32

Let Δ := (t, s) ∈ R2 : a < s < t < b .

330

Appendix A

(a) Let k ∈ L2 (Δ). Then the Volterra operator

t k(t, s) x(s) ds ,

(Kx)(t) :=

a < t < b , x ∈ L2 (a, b) ,

(A.20)

a

is bounded in L2 (a, b). There exists m ∈ N with K m L(L2 (a,b)) < 1. The Volterra equation of the second kind

t x(t) −

k(t, s) x(s) ds = y(t) ,

a < t < b,

(A.21)

a

is uniquely solvable in L2 (a, b) for every y ∈ L2 (a, b), and the solution x depends continuously on y. The solution x ∈ L2 (a, b) has the form

t r(t, s) y(s) ds ,

x(t) = y(t) +

t ∈ (a, b) ,

a

with some kernel r ∈ L2 (Δ). (b) Let k ∈ C(Δ). Then the operator K defined by (A.20) is bounded in C[a, b], and there exists m ∈ N with K m ∞ < 1. Equation (A.21) is also uniquely solvable in C[a, b] for every y ∈ C[a, b], and the solution x depends continuously on y. For the remaining part of this section, we assume that X and Y are normed spaces and K : X → Y a linear and bounded operator. Definition A.33 (Compact Operator) The operator K : X → Y is called compact if it maps every bounded set S into a relatively compact set K(S). We recall that a set M ⊂ Y is called relatively compact if every bounded sequence (yj )j in M has an accumulation point in closure(M ); that is, if the closure closure(M ) is compact. The set of all compact operators from X into Y is a closed subspace of L(X, Y ) and even a two-sided ideal by part (c) of the following theorem. Theorem A.34 (a) If K1 and K2 are compact from X into Y , then so are K1 + K2 and λK1 for every λ ∈ K. (b) Let Kn : X −→ Y be a sequence of compact operators between Banach spaces X and Y . Let K : X −→ Y be bounded, and let Kn converge to K in the operator norm; that is, Kn − KL(X,Y ) := sup x=0

Kn x − KxY −→ 0 (n −→ ∞) . xX

A.3 Linear Bounded and Compact Operators

331

Then K is also compact. (c) If L ∈ L(X, Y ) and K ∈ L(Y, Z), and L or K is compact, then KL is also compact. (d) Let An ∈ L(X, Y ) be pointwise convergent to some A ∈ L(X, Y ); that is, An x → Ax for all x ∈ X. If K : Z → X is compact, then An K−A KL(Z.Y ) → 0; that is, the operators An K converge to AK in the operator norm. (e) The identity operator x → x is compact as an operator from X into itself if, and only if, X is finite dimensional. (f ) Every bounded operator K from X into Y with finite-dimensional range is compact.   Theorem A.35 (a) Let k ∈ L2 (c, d) × (a, b) . The operator K : L2 (a, b) → L2 (c, d), defined by

b (Kx)(t) :=

k(t, s) x(s) ds ,

t ∈ (c, d) ,

x ∈ L2 (a, b) ,

(A.22)

a

is compact from L2 (a, b) into L2 (c, d). (b) Let k be continuous on [c, d]×[a, b] or weakly singular on [a, b]×[a, b] (in this case [c, d] = [a, b]). Then K defined by (A.22) is also compact as an operator from C[a, b] into C[c, d]. We now study equations of the form x − Kx = y ,

(A.23)

where the linear operator K : X → X is compact. The following theorem extends the well-known existence results for finite linear systems of n equations and n variables to compact perturbations of the identity. Theorem A.36 (Riesz) Let X be a normed space and K : X → X be a linear compact operator. (a) The null space N (I − K) = {x ∈ X : x = Kx} is finite-dimensional and the range R(I − K) = (I − K)(X) is closed in X. (b) If I − K is one-to-one, then I − K is also surjective, and the inverse (I − K)−1 is bounded. In other words, if the homogeneous equation x − Kx = 0 admits only the trivial solution x = 0, then the inhomogeneous equation x − Kx = y is uniquely solvable for every y ∈ X and the solution x depends continuously on y. The next theorem studies approximations of equations of the form Ax = y. Again, we have in mind that A = I − K.

332

Appendix A

Theorem A.37 Assume that the operator A : X → Y between Banach spaces X and Y has a bounded inverse A−1 . Let An ∈ L(X, Y ) be a sequence of bounded operators that converge in norm to A; that is, An − AL(X,Y ) → 0 as n → ∞. Then, for sufficiently large n, more precisely for all n with A−1 (An − A)L(X) < 1 ,

(A.24)

the inverse operators A−1 n : Y → X exist and are uniformly bounded by A−1 n L(Y,X) ≤

A−1 L(Y,X) ≤ c. 1 − A−1 (An − A)L(X)

(A.25)

For the solutions of the equations Ax = y

and

An xn = yn ,

the error estimate

xn − xX ≤ c An x − AxY + yn − yY

(A.26)

holds with the constant c from (A.25).

A.4

Sobolev Spaces of Periodic Functions

In this section, we recall definitions and properties of Sobolev (Hilbert) spaces of periodic functions. A complete discussion including proofs can be found in the monograph [168]. From Parseval’s identity (A.11), we note that x ∈ L2 (0, 2π) if and only if the Fourier coefficients

2π 1 ak = x(s) e−iks ds , k ∈ Z , (A.27) 2π 0

are square summable. In this case  |ak |2 = k∈Z

1 x2L2 . 2π

If x is periodic and continuously differentiable on [0, 2π], partial integration of (A.27) yields the formula ak

−i = 2πk

2π

x (s) e−iks ds ;

0

that is, ikak are the Fourier coefficients of x and are thus square summable. r (0, 2π) of L2 (0, 2π) by requirThis motivates the introduction of subspaces Hper ing for their elements a certain decay of the Fourier coefficients ak . In the following we set ψˆk (t) = eikt

for k ∈ Z and t ∈ (0, 2π) .

A.4 Sobolev Spaces of Periodic Functions

333

Definition A.38 (Sobolev space of periodic functions) r For r ≥ 0, the Sobolev space Hper (0, 2π) of order r is defined by     r Hper (0, 2π) := x = ak ψˆk ∈ L2 (0, 2π) : (1 + k 2 )r |ak |2 < ∞ . k∈Z 0 (0, 2π) Hper

We note that

k∈Z 2

coincides with L (0, 2π).

r Theorem A.39 The Sobolev space Hper (0, 2π) is a Hilbert space with the inner product defined by  r (x, y)Hper := (1 + k 2 )r ak bk , (A.28)

where x = by

k∈Z



 r ˆ ˆ k∈Z ak ψk and y = k∈Z bk ψk . The norm in Hper (0, 2π) is given  r = xHper



1/2 (1 + k 2 )r |ak |2

.

k∈Z r The Sobolev space Hper (0, 2π) is dense in L2 (0, 2π). √ 0 0 We note that xL2 = 2π xHper , that is, the norms xL2 and xHper are equivalent on L2 (0, 2π). r [0, 2π] = x ∈ Theorem A.40 (a) For r ∈ N0 := N ∪ {0}, the space Cper

r (0, 2π). C r [0, 2π] : x is 2π − periodic is boundedly embedded in Hper (b) The space T of all trigonometric polynomials   n T := ak ψˆk : ak ∈ C, n ∈ N k=−n

is dense in

r Hper (0, 2π)

for every r ≥ 0.

⊂ L2 (0, 2π) as a Banach space (that is, forget We consider r (0, 2π) → the inner product for a moment) with bounded imbedding j : Hper 2 L (0, 2π) which has a dense range. Therefore, we can consider the correspondr r r (0, 2π) ⊂ L2 (0, 2π) ⊂ Hper (0, 2π) where Hper (0, 2π) ing Gelfand triple Hper r denotes the space of all anti-linear functionals on Hper (0, 2π), see Definition A.26. We make the following definition. r (0, 2π) Hper

−r r (0, 2π) = Hper (0, 2π) the antiDefinition A.41 For r ≥ 0, we denote by Hper r (0, 2π); that is, the space of all anti-linear bounded functionals dual space of Hper r r −r (0, 2π). Then Hper (0, 2π) ⊂ L2 (0, 2π) ⊂ Hper (0, 2π) with bounded and on Hper −r (0, 2π) × dense imbeddings. The corresponding sesquilinear form ·, · : Hper r (0, 2π) → C extends the inner product in L2 (0, 2π); that is, Hper

2π ψ, φ = (ψ, φ)

L2

=

ψ(t) φ(t) dt 0

r for all ψ ∈ L2 (0, 2π) and φ ∈ Hper (0, 2π).

334

Appendix A

The following theorems give characterizations in terms of the Fourier coefficients. Theorem A.42 Let again ψˆk (t) = eikt for k ∈ Z and t ∈ (0, 2π). −r r (a) Let  ∈ Hper (0, 2π) = Hper (0, 2π) and define ck := , ψˆk for k ∈ Z. Then  1/2 −r Hper = (1 + k 2 )−r |ck |2 k∈Z

and , x =



ck ak

for all x =

k∈Z



r ak ψˆk ∈ Hper (0, 2π) .

(A.29)

k∈Z

(b) Conversely, let ck ∈ C satisfy  (1 + k 2 )−r |ck |2 < ∞ . k∈Z −r Then , defined by (A.29), is in Hper (0, 2π) with 2 H −r .



k∈Z (1

+ k 2 )−r |ck |2 =

per

 Proof: (a) Set z N = |k|≤N ck (1 + k 2 )−r ψˆk for N ∈ N. Then (note that  is anti-linear)  N r −r z (1 + k 2 )−r |ck |2 = , z N ≤ Hper Hper |k|≤N −r = Hper

 

(1 + k 2 )−2r |ck |2 (1 + k 2 )r

|k|≤N

−r = Hper

 

(1 + k 2 )−r |ck |2

|k|≤N



+ k 2 )−r |ck |2 ≤ 2H −r by letting N tend to infinity. In per  r particular, the series converges. Furthermore, for x = k∈Z ak ψˆk ∈ Hper (0, 2π)   N N we set x = |k|≤N ak ψˆk , and have that , x = |k|≤N ak ck and thus

which proves

  , xN 

k∈Z (1

      =  ak ck  ≤ |k|≤N

≤

 

 |k|≤N

|ak |2 (1 + k 2 )r

|k|≤N r = xN Hper

  |k|≤N

|ak |(1 + k 2 )r/2 |ck |(1 + k 2 )−r/2  

|ck |2 (1 + k 2 )−r

|k|≤N

|ck |2 (1 + k 2 )−r

A.4 Sobolev Spaces of Periodic Functions and thus also , x = N tend to infinity.

 k∈Z

335

ak ck and 2H −r ≤



per

k∈Z (1+k

2 −r

)

|ck |2 by letting



(b) This is shown in the same way.

−r This theorem makes it possible to equip the Banach space Hper (0, 2π) = with an inner product and make it a Hilbert space.

r (0, 2π) Hper

Theorem A.43 Let r > 0. On L2 (0, 2π), we define the inner product and norm by  (x, y)−r := (1 + k 2 )−r ak bk , (A.30a) k∈Z

x−r

:=



(1 + k 2 )−r |ak |2 ,

(A.30b)

k∈Z

  respectively, where x = k∈Z ak ψˆk and y = k∈Z bk ψˆk . Then the completion ˜ −r (0, 2π) of L2 (0, 2π) with respect to  · −r can be identified with H −r (0, 2π). H per per −r (0, 2π), where The isomorphism is given by the extension of J : L2 (0, 2π) → Hper Jx, y :=



ak bk

for x =

k∈Z



r ak ψˆk ∈ Hper (0, 2π)

k∈Z

 −r and y = k∈Z bk ψˆk ∈ L2 (0, 2π). Therefore, we identify p−r with JpHper −r . and simply write pHper −r Proof: First we show that Jx ∈ Hper (0, 2π). Indeed, by the Cauchy–Schwarz  r (0, 2π) inequality, we have for y = k∈Z bk ψˆk ∈ Hper

|Jx, y |

≤



(1 + k 2 )−r/2 |ak |

k∈Z

 ≤





(1 + k 2 )r/2 |bk |

1/2 ⎛ ⎝

(1 + k 2 )−r |ak |2

⎞1/2 (1 + k 2 )r |bk |2 ⎠

|k|≤N

k∈Z

≤



−r yr , xHper

−r 2 −r ≤ x−r for all x ∈ L (0, 2π). By the (0, 2π) with JxHper and thus Jx ∈ Hper  previous theorem, applied to  = Jx, we have Jx2H −r = k∈Z (1 + k 2 )−r |ck |2 per with ck = Jx, ψˆk = ak . Therefore, Jx −r = x−r , and J can be extended Hper

˜ −r (0, 2π) into H −r (0, 2π) (by Theorem A.30). to a bounded operator from H per per −r r (0, 2π) = Hper (0, 2π) and It remains to show that J is surjective. Let  ∈ Hper   define ck = , ψˆk and xN = |k|≤N ck ψˆk for N ∈ N. For y = k∈Z bk ψˆk ∈     r Hper (0, 2π) we have JxN = |k|≤N bk ck = , |k|≤N bk ψˆk which converges

336

Appendix A

to , y . Furthermore, (xN ) is a Cauchy sequence with respect to ·−r because of the convergence of k∈Z (1 + k 2 )−r |ck |2 by the previous theorem, part (a).  r Theorem A.44 (a) For r > s, the Sobolev space Hper (0, 2π) is a dense subspace s r s (0, 2π) is of Hper (0, 2π). The inclusion operator from Hper (0, 2π) into Hper compact. r (b) For all r ≥ 0 and x ∈ L2 (0, 2π) and y ∈ Hper (0, 2π), there holds     (x, y)L2  = 2π (x, y)H 0  ≤ 2π x −r yH r . Hper per per

(A.31)

We note that the estimate (A.31) is in accordance with (A.16) because, with −r the imbedding j  from L2 (0, 2π) into Hper (0, 2π) and the identification J from −r −r ˜ per (0, 2π) onto Hper (0, 2π) of the previous theorem we show easily that the H −r ˜ per imbedding of L2 (0, 2π) into H (0, 2π) is given by J −1 ◦ j  where (J −1 ◦ j  )x = 2 2πx for x ∈ L (0, 2π). Proof: (a) Denseness is easily seen by truncating the Fourier series. Compactness is shown by defining the finite-dimensional operators  JN from r s ˆ ˆ (0, 2π) into Hper (0, 2π) by JN x = a where x = ψ Hper |k|≤N k k k∈Z ak ψk . Then JN is compact by part (f) of Theorem A.34 and JN x − x2Hper s

=



(1 + k 2 )s |ak |2 ≤

|k|>N

≤

1 N 2(r−s)

 1 (1 + k 2 )r |ak |2 (1 + N 2 )r−s |k|>N

x2Hper . r

Therefore, JN converges in the operator norm to the imbedding J which implies, again by Theorem A.34, that also J is compact. (b) This is seen as in the proof of Theorem A.43 with the Cauchy–Schwarz inequality.  Theorems A.40 and A.43 imply that the space T of all trigonometric polyr nomials is dense in Hper (0, 2π) for every r ∈ R. Now we study the orthogonal projection and the interpolation operators with respect to equidistant knots and the 2n-dimensional space  n−1   Tn := ak ψˆk : ak ∈ C (A.32) k=−n

where again ψˆk (t) = eikt for t ∈ [0, 2π] and k ∈ Z. Lemma A.45 Let r, s ∈ R with r ≥ s. (a) The following stability estimate holds r s zn Hper ≤ c nr−s zn Hper

for all zn ∈ Tn .

A.4 Sobolev Spaces of Periodic Functions

337

(b) Let Pn : L2 (0, 2π) → Tn ⊂ L2 (0, 2π) be the orthogonal projection operator. Then Pn is given by n−1 

Pn x =

ak ψˆk ,

x ∈ L2 (0, 2π) ,

(A.33)

k=−n

where 1 = 2π

ak

2π

x(s) e−iks ds ,

k ∈ Z,

0

are the Fourier coefficients of x. Furthermore, the following estimate holds: s ≤ x − Pn xHper

Proof:

(a) With zn =

zn 2Hper r

≤



1 r xHper nr−s

n−1 k=−n

(A.34)

ak ψˆk this follows simply by

(1 + k 2 )r |ak |2 =

|k|≤n

≤

r for all x ∈ Hper (0, 2π) ,



(1 + k 2 )r−s (1 + k 2 )s |ak |2

|k|≤n

(1 + n2 )r−s zn 2Hper ≤ (2n2 )r−s zn 2Hper . s s

 (b) Let x = k∈Z ak ψˆk ∈ L2 (0, 2π) and define the right-hand side of (A.33) n−1 by z; that is, z = k=−n ak ψˆk ∈ Tn . The orthogonality of ψˆk implies that x − z is orthogonal to Tn . This proves that z coincides with Pn x. Now let r x ∈ Hper (0, 2π). Then x − Pn x2Hper s

≤



(1 + k 2 )s |ak |2

|k|≥n

=



(1 + k 2 )−(r−s) (1 + k 2 )r |ak |2

!

|k|≥n

≤

(1 + n2 )s−r x2Hper ≤ n2(s−r) x2Hper . r r



Now let tj := j nπ , j = 0, . . . , 2n − 1, be equidistantly chosen points in [0, 2π]. Interpolation of smooth periodic functions by trigonometric polynomials can be found in numerous books as, for example, in [72]. Interpolation in Sobolev spaces of integer orders can be found in [42]. We give a different and much simpler proof of the error estimates that are optimal and hold in Sobolev spaces of fractional order. Theorem A.46 For every n ∈ N and every 2π-periodic function x ∈ C[0, 2π], there exists a unique pn ∈ Tn with x(tj ) = pn (tj ) for all j = 0, . . . , 2n − 1. The

338

Appendix A

trigonometric interpolation operator Qn : Cper [0, 2π] = x ∈ C[0, 2π] : x(0) =

x(2π) → Tn has the form Qn x =

2n−1 

x(tk ) Lk

k=0

with Lagrange interpolation basis functions Lk (t) =

n−1 1  im(t−tk ) e , 2n m=−n

k = 0, . . . , 2n − 1 .

(A.35)

The interpolation operator Qn has an extension to a bounded operator from r r (0, 2π) into Tn ⊂ Hper (0, 2π) for all r > 12 . Furthermore, Qn obeys estiHper mates of the form s x − Qn xHper ≤

c r xHper nr−s

r for all x ∈ Hper (0, 2π) ,

(A.36)

where 0 ≤ s ≤ r and r > 12 . The constant c depends only on s and r. In r (0,2π)) is uniformly bounded with respect to n. particular, Qn L(Hper Proof:  The proof of the first part can be found in, for example, [168]. Let x(t) = m∈Z am exp(imt). Direct calculation shows that for smooth functions x the interpolation is given by (Qn x)(t) =

n−1 

a ˆj eijt

with

j=−n

a ˆj

=

2n−1 1  x(tk ) e−ijkπ/n , 2n

j = −n, . . . , n − 1 .

k=0

The connection between the continuous and discrete Fourier coefficients is simply a ˆj

=

2n−1 1   am eimkπ/n−ijkπ/n 2n

=

2n−1 #k  "  1  ei(m−j)π/n am = aj+2n 2n

k=0 m∈Z

m∈Z

k=0

∈Z



where x = m∈Z am ψˆm . It is sufficient to estimate Pn x − Qn x because the required estimate holds for x − Pn x by formula (A.34). We have Pn x − Qn x =

n−1 

[am − a ˆ m ] ψm

m=−n

A.4 Sobolev Spaces of Periodic Functions

339

and thus by the Cauchy–Schwarz inequality Pn x − Qn x2Hper s =

n−1 

|am − a ˆm |2 (1 + m2 )s ≤ c n2s

≤

 2     a m+2n  

m=−n =0

m=−n

≤

n−1 

2  r/2  1 + (m + 2n)2 m=−n =0 % n−1  $  r 1  r . 1 + (m + 2n)2 |am+2n |2 c n2s 1 + (m + 2n)2 m=−n =0 =0 c n2s

n−1 

   r/2 !  1 + (m + 2n)2 am+2n  

1

From the obvious estimate '−2r & m  −r + 1 + (m + 2n)2 ≤ (2n)−2r ≤ c n−2r 2n =0

=0

for all |m| ≤ n and n ∈ N, we conclude that Pn x − Qn x2Hper s

≤

c n2(s−r)

≤

2(s−r)

n−1 

 r 1 + (m + 2n)2 |am+2n |2

m=−n =0

cn

x2Hper . r



For real-valued functions, it is more convenient to study the orthogonal projection and interpolation in the 2n-dimensional space ⎧ ⎫ n n−1 ⎨ ⎬  aj cos(jt) + bj sin(jt) : aj , bj ∈ R . ⎩ ⎭ j=0

j=1

In this case, the Lagrange interpolation basis functions Lk are given by (see [168])   n−1  1 1 + 2 cos m(t − tk ) + cos n(t − tk ) , (A.37) Lk (t) = 2n m=1 k = 0, . . . , 2n − 1, and the estimates (A.34) and (A.36) are proven by the same arguments.   Theorem A.47 Let r ∈ N and k ∈ C r [0, 2π] × [0, 2π] be 2π-periodic with respect to both variables. Then the integral operator K, defined by

2π (Kx)(t) :=

k(t, s) x(s) ds ,

t ∈ (0, 2π) ,

(A.38)

0 p r can be extended to a bounded operator from Hper (0, 2π) into Hper (0, 2π) for every −r ≤ p ≤ r.

340 Proof:

Appendix A Let x ∈ L2 (0, 2π). From dj (Kx)(t) = dtj

2π 0

∂ j k(t, s) x(s) ds , ∂tj

j = 0, . . . , r ,

we conclude from Theorem A.44 that for x ∈ L2 (0, 2π)  j   j   ∂ k(t, ·)  d      −r  dtj (Kx)(t) ≤ 2π  ∂tj  r xHper H per

and thus r −r ≤ c1 KxC r ≤ c2 xHper KxHper

for all x ∈ L2 (0, 2π). Application of Theorem A.30 yields the assertion because −r (0, 2π).  L2 (0, 2π) is dense in Hper

A.5

Sobolev Spaces on the Unit Disc

Let B = {x ∈ R2 : |x| < 1} be the open unit disc with boundary ∂B. In this section, we consider functions from B into C that we describe by Cartesian coordinates x = (x1 , x2 ) or by polar coordinates (r, ϕ). Functions on the boundary are identified with 2π-periodic functions on R. As in the case of the s (0, 2π) we define the Sobolev space H 1 (B) by completion. Sobolev spaces Hper Definition A.48 The Sobolev space H 1 (B) is defined as the completion of C ∞ (B) with respect to the norm 

 !  f H 1 (B) =  |f (x)|2 + |∇f (x)|2 dx . (A.39) B

We express the norm (A.39) in polar coordinates (r, ϕ). The gradient is given in polar coordinates as 1 ∂f (r, ϕ) ∂f (r, ϕ) rˆ + ϕˆ , ∂r r ∂ϕ  ϕ  sin ϕ where rˆ = cos ˆ = −cos denote the unit vectors. We fix r > 0 and sin ϕ and ϕ ϕ expand the function f (r, ·) (formally) into a Fourier series with respect to ϕ:  f (r, ϕ) = fm (r) eimϕ ∇f (r, ϕ) =

m∈Z

with Fourier coefficients 1 fm (r) = 2π

2π 0

f (r, t) e−imt dt ,

m ∈ Z,

A.5 Sobolev Spaces on the Unit Disc

341

that depend on r. Therefore,  ∂f (r, ϕ)  = fm (r) eimϕ , ∂r m∈Z

 ∂f (r, ϕ) = i fm (r) m eimϕ . ∂ϕ m∈Z

The norm in H 1 (B) is given by f 2H 1 (B)

1  %  $ m2  = 2π (r)|2 r dr , 1 + 2 |fm (r)|2 + |fm r

(A.40)

m∈Z 0

because |∇f (r, ϕ)|

2

   ∂f (r, ϕ) 2  + 1  =  ∂r  r2

   ∂f (r, ϕ) 2    ∂ϕ 

and 2

2π     imϕ   f (r) e |fm (r)|2 . m   dϕ = 2π 0

m∈Z

m∈Z

To every function f ∈ C ∞ (B), one can assign the trace f |∂B on ∂B. We denote this mapping by τ , thus τ : C ∞ (B) → C ∞ (∂B) is defined as τ f = f |∂B . The following result is central. Theorem A.49 (Trace Theorem) The trace operator τ has an extension to a bounded operator from H 1 (B) to 1/2 H 1/2 (∂B), where again H 1/2 (∂B) is identified with the Sobolev space Hper (0, 2π) 1 of periodic functions (see Definition A.38). Furthermore, τ : H (B)→H 1/2 (∂B) is surjective. More precisely, there exists a bounded linear operator E : H 1/2 (∂B) → H 1 (B) with τ ◦ E = I on H 1/2 (∂B) (that is, E is a right inverse of τ ). Proof: Let f ∈ C ∞ (B). Then (see (A.40) and (A.28)) f 2H 1 (B)

=

τ f 2H 1/2 (∂B)

=

1  %  $ m2 2  2 2π 1 + 2 |fm (r)| + |fm (r)| r dr , r m∈Z 0  1 + m2 |fm (1)|2 . m∈Z

We estimate, using the fundamental theorem of calculus and the inequality of Cauchy–Schwarz,

342

|fm (1)|

Appendix A

2

1 = 0

 d 2 r |fm (r)|2 dr dr

1 =

1

2

|fm (r)| r dr + 2 Re

2 0

≤

 (r) r 2 dr fm (r) fm

0

   1  1

1      (r)|2 r 2 dr . 2 |fm (r)|2 r dr + 2  |fm (r)|2 r2 dr |fm 0

0

Using the inequality 2ab ≤  1 + m2 |fm (1)|2

√

≤ 2

2 √ b 1+m2

1 + m2 a2 +



1 1+

0

yields

|fm (r)|2 r dr

m2 0 2

1

|fm (r)| r dr +

+ (1 + m ) 0

≤ 3 (1 + m2 )

1  ≤ 3 0

1

2 2

1 0 2

m 1+ 2 r

|fm (r)|2 r dr +

1 0



2

 |fm (r)|2 r2 dr

0  |fm (r)|2 r dr

1

|fm (r)| r dr +

 |fm (r)|2 r dr ,

0

where we have also used the estimates r2 ≤ r for r ∈ [0, 1] and 1 + m2 . By summation, we conclude that τ f 2H 1/2 (∂B) ≤

3 f 2H 1 (B) 2π

√ 1 + m2 ≤

for all f ∈ C ∞ (B) .

(A.41)

Therefore, the trace operator is bounded with respect to the norms of H 1 (B) and H 1/2 (∂B). By the general functional analytic Theorem A.30, the operator has an extension to a bounded operator from H 1 (B) into H 1/2 (∂B). We define the operator E : C ∞ (∂B) → C ∞ (B) by (Ef )(r, ϕ) =



fm r|m| eimϕ ,

r ∈ [0, 1] , ϕ ∈ [0, 2π] .

m∈Z

Here, again, fm are the Fourier coefficients of f ∈ C ∞ (∂B).  Obviously, (τ Ef )(ϕ) = m∈Z fm eimϕ = f (ϕ); that is, E is a right inverse of τ .

A.5 Sobolev Spaces on the Unit Disc

343

It remains to show the boundedness of E. 1  %  $ m2 2 2 2|m| 2 2 2|m|−2 Ef H 1 (B) = 2π + |fm | m r 1 + 2 |fm | r r dr , r m∈Z 0     1 + |m| ≤ 2π = 2π |fm |2 |fm |2 (1 + |m|) 2|m| + 2 m∈Z m∈Z   √ √ 2 ≤ 2 2π |fm | 1 + m2 = 2 2πf 2H 1/2 (∂B) m∈Z

√ √ where we used the inequality 1 + |m| ≤ 2 1 + m2 . Therefore, E also pos sesses an extension to a bounded operator from H 1/2 (∂B) to H 1 (B). Remark: The trace operator is compact when considered as an operator from H 1 (B) to L2 (∂B) because it is the composition of the bounded operator τ : H 1 (B) → H 1/2 (∂B) and the compact embedding j : H 1/2 (∂B) → L2 (∂B). We now consider the subspaces  

f ∈ L2 (∂B) : f d = 0 , L2 (∂B) = ∂B  

1/2 f ∈ H 1/2 (∂B) : f d = 0 , H (∂B) = ∂B  

f ∈ H 1 (B) : τ f d = 0 . H 1 (B) = ∂B

 2π

1/2

Because 0 exp(imϕ) dϕ = 0 for m = 0, the spaces H (∂B) and H 1 (B) consist exactly of the functions with the representations  fm eimϕ and f (ϕ) = m∈Z m=0

f (r, ϕ)

=



fm (r) eimϕ ,

m∈Z

that satisfy the summation conditions  1 + m2 |fm |2 < ∞

and

m∈Z m=0 1  %  $ m2  (r)|2 r dr < ∞ 1 + 2 |fm (r)|2 + |fm r

and f0 (1) = 0 ,

m∈Z 0

respectively. We can define an equivalent norm in the subspace H 1 (B). This is a consequence of the following result.

344

Appendix A

Theorem A.50 (Friedrich’s Inequality) For all f ∈ H 1 (B), we have √ f L2 (B) ≤ 2 ∇f L2 (B) .

(A.42)

Proof: Again, we use the representation of the norm in polar coordinates: r |fm (r)|

2

r = 0

 d s |fm (s)|2 ds ds

r

r

2

0

1 ≤

0

   1  1     2 2  (s)|2 s ds  |fm (s)| ds + 2 |fm (s)| s ds |fm

0

1 ≤

 (s) s ds fm (s) fm

|fm (s)| ds + 2 Re

=

0

0

1

|fm (s)|2 (1 + s) ds +

0

 |fm (s)|2 s ds ,

0 2

2

where we again used 2ab ≤ a + b in the last step. First let |m| ≥ 1. By 1 + s ≤ 2 ≤ 2m2 /s, it is r |fm (r)|2 ≤ 2

1 0

m2 |fm (s)|2 s ds + s2

1

 |fm (s)|2 s ds ,

0

and thus by integration

1

1 

2

r |fm (r)| dr ≤ 2 0

0

 m2 2  2 |f (s)| + |f (s)| s ds . m m s2

We finally consider f0 . It is f0 (r) = − r |f0 (r)|

2

1 ≤ (1 − r) r

|f0 (s)|2 r ds

1 r

(A.43)

f0 (s) ds because f0 (1) = 0, thus

1

≤ r

|f0 (s)|2 s ds

1 ≤

|f0 (s)|2 s ds .

0

Therefore, (A.43) also holds for m = 0. Summation with respect to m yields the assertion.  Remark: Therefore, f → ∇f L2 (B) defines an equivalent norm to  · H 1 (B) in H 1 (B). Indeed, for f ∈ H 1 (B) it holds by Friedrich’s inequality: f 2H 1 (B) = f 2L2 (B) + ∇f 2L2 (B) ≤ 3 ∇f 2L2 (B) ,

A.6 Spectral Theory for Compact Operators in Hilbert Spaces

345

thus 1 √ f H 1 (B) ≤ ∇f L2 (B) ≤ f H 1 (B) 3

for all f ∈ H 1 (B) .

(A.44)

So far, we considered spaces of complex-valued functions. The spaces of realvalued functions are closed subspaces. In the Fourier representation, one has    f (r, ϕ) = fm (r) e−imϕ = f−m (r) eimϕ = f (r, ϕ) = fm (r) eimϕ m∈Z

m∈Z

m∈Z

because f (r, ϕ) = f (r, ϕ). Therefore, f−m = fm for all m. All of the theorems remain valid also for Sobolev spaces of real-valued functions.

A.6

Spectral Theory for Compact Operators in Hilbert Spaces

Definition A.51 (Spectrum) Let X be a normed space and A : X −→ X be a linear operator. The spectrum σ(A) is defined as the set of (complex) numbers λ such that the operator A − λI does not have a bounded inverse on X. Here, I denotes the identity on X. λ ∈ σ(A) is called an eigenvalue of A if A − λI is not one-to-one. If λ is an eigenvalue, then the nontrivial elements x of the kernel N (A − λI) = {x ∈ X : Ax − λx = 0} are called eigenvectors of A. This definition makes sense for arbitrary linear operators in normed spaces. For noncompact operators A, it is possible that even for λ = 0 the operator A − λI is one-to-one but fails to be bijective. As an example, we consider X = 2 and define A by  0, if k = 1 , (Ax)k := xk−1 , if k ≥ 2 , for x = (xk ) ∈ 2 . Then λ = 1 belongs to the spectrum of A but is not an eigenvalue of A, see Problem A.4. Theorem A.52 Let A : X → X be a linear bounded operator. (a) Let xj ∈ X, j = 1, . . . , n, be a finite set of eigenvectors corresponding to pairwise different eigenvalues λj ∈ C. Then {x1 , . . . , xn } are linearly independent. If X is a Hilbert space and A is self-adjoint (that is, A∗ = A), then all eigenvalues λj are real-valued and the corresponding eigenvectors x1 , . . . , xn are pairwise orthogonal. (b) Let X be a Hilbert space and A : X → X be self-adjoint. Then AL(X) =

sup |(Ax, x)X | = r(A) ,

x X =1

where r(A) = sup{|λ| : λ ∈ σ(A)} is called the spectral radius of A.

346

Appendix A

The situation is simpler for compact operators. We collect the most important results in the following fundamental theorem. Theorem A.53 (Spectral Theorem for Compact Self-Adjoint Operators) Let K : X → X be compact and self-adjoint (and K = 0). Then the following holds: (a) The spectrum consists only of eigenvalues and possibly 0. Every eigenvalue of K is real-valued. K has at least one but at most a countable number of eigenvalues with 0 as the only possible accumulation point. (b) For every eigenvalue λ = 0, there exist only finitely many linearly independent eigenvectors; that is, the eigenspaces are finite-dimensional. Eigenvectors corresponding to different eigenvalues are orthogonal. (c) We order the eigenvalues in the form |λ1 | ≥ |λ2 | ≥ |λ3 | ≥ . . . and denote by Pj : X → N (K − λj I) the orthogonal projection onto the eigenspace corresponding to λj . If there exist only a finite number λ1 , . . . , λm of eigenvalues, then K =

m 

λ j Pj .

j=1

If there exists an infinite sequence (λj ) of eigenvalues, then K =

∞ 

λ j Pj ,

j=1

where the series converges in the operator norm. Furthermore,   m     K − λ P = |λm+1 | . j j  j=1

L(X)

(d) Let H be the linear span of all of the eigenvectors corresponding to the eigenvalues λj = 0 of K. Then X = closure(H) ⊕ N (K) . Sometimes, part (d) is formulated differently. For a common treatment of the cases of finitely and infinitely many eigenvalues, we introduce the index set J ⊂ N, where J is finite in the first case and J = N in the second case. For every eigenvalue λj , j ∈ J, we choose an orthonormal basis of the corresponding eigenspace N (K − λj I). Again, let the eigenvalues λj = 0 be ordered in the form |λ1 | ≥ |λ2 | ≥ |λ3 | ≥ . . . > 0 .

A.6 Spectral Theory for Compact Operators in Hilbert Spaces

347

By counting every λj = 0 relative to its multiplicity, we can assign an eigenvector xj to every eigenvalue λj . Then every x ∈ X possesses an abstract Fourier expansion of the form x = x0 +



(x, xj )X xj

j∈J

for some x0 ∈ N (K) and Kx =



λj (x, xj )X xj .

j∈J

As a corollary, we observe that the set {xj : j ∈ J} of all eigenvectors forms a complete system in X if K is one-to-one. The eigenvalues can be expressed by Courant’s max-inf and min-sup principle. We need it in the following form. Theorem A.54 Let K : X → X be compact and self-adjoint (and K = 0) and + let {λ− j : j = 1, . . . , n− } and {λj : j = 1, . . . , n+ } be its negative and positive eigenvalues, respectively, ordered as − − + + + λ− 1 ≤ λ 2 ≤ · · · ≤ λ j ≤ · · · < 0 < · · · ≤ λj ≤ · · · ≤ λ 2 ≤ λ 1

and counted according to multiplicity. Here n± can be zero (if no positive or negative, respectively, eigenvalues occur), finite, or infinity. (a) If there exist positive eigenvalues λ+ m then λ+ m =

min

sup

V ⊂X x∈V ⊥ dim V =m−1

(Kx, x)X . x2X

(b) If there exist negative eigenvalues λ− m then λ− m =

max

inf

V ⊂X x∈V dim V =m−1

⊥

(Kx, x)X . x2X

Proof: We only prove part (b) because this is needed in the proof of Theorem 7.52. Set (Kx, x)X μm = sup inf ⊥ x2X x∈V V ⊂X dim V =m−1

for any x ∈ X we use the representation as x = x0 +  + For  abbreviation. − − + c x + c x with Kx0 = 0. Here, x± j are the eigenvectors correspondj j j j j j    − −2 ± − 2 + + 2 ing to λj . Then (Kx, x)X = j λ− |c | + j j j λj |cj | ≥ j λj |cj | .

348

Appendix A

⊥ For V = span{x− we have that c− j : j = 1, . . . , m − 1} and x ∈ V j = − (x, xj )X = 0 for j = 1, . . . , m − 1 and thus

(Kx, x)X ≥ inf ⊥ x2X x∈V

 inf

x∈V ⊥

− 2 λ− j |cj |  + 2 ≥ λ− m − 2 |c | + j j j |cj |

j≥m

x0 2X +



− because λ− j < 0. Therefore, μm ≥ λm . We note that for this choice of V the − infimum is attained for x = xm . Let, on the other hand, V ⊂ X be an arbitrary subspace of dimension m − 1. ˆ ∈ span{x− Choose a basis {vj : j = 1, . . . , m − 1} of V . Construct x j : j = 1, . . . , m} such that x ˆ ⊥ V and ˆ xX= 1. This is possible. Indeed, the ansatz m m − x ˆ = j=1 cj x− j leads to the system j=1 cj (xj , v )X = 0 for  = 1, . . . , m − 1. This system of m − 1 equations and mvariables has a nontrivial solution which m can be normalized such that ˆ x2X = j=1 |cj |2 = 1. Therefore,

inf

x∈V ⊥

m m   (Kx, x)X − 2 − ≤ (K x ˆ , x ˆ ) = λ |c | ≤ λ |cj |2 = λ− X j m m. j x2X j=1 j=1

This shows μm ≤ λ− m and ends the proof.



The following corollary is helpful. Corollary A.55 Let K : X → X and λ± j as in the previous theorem. (a) If there exists a subspace W ⊂ X of dimension m with (Kx, x)X +x2X ≤ 0 for all x ∈ W then (at least) m negative eigenvalues exist and λ− 1 ≤ ≤ −1. · · · 0 the ordered sequence of the positive singular values of K, counted relative to its multiplicity. Then there exist orthonormal systems {xj : j ∈ J} ⊂ X and {yj : j ∈ J} ⊂ Y with the following properties: Kxj = μj yj

and

K ∗ yj = μj xj

for all j ∈ J .

The system {μj , xj , yj : j ∈ J} is called a singular system for K. Every x ∈ X possesses the singular value decomposition  x = x0 + (x, xj )X xj j∈J

for some x0 ∈ N (K) and Kx =



μj (x, xj )X yj .

j∈J

Note that J ⊂ N can be finite or infinite; that is, J = N.

350

Appendix A

The following theorem characterizes the range of a compact operator with the help of a singular system. Theorem A.58 (Picard) Let K : X −→ Y be a linear compact operator with singular system {μj , xj , yj : j ∈ J}. The equation Kx = y (A.45) is solvable if and only if y ∈ N (K ∗ )⊥

 1 |(y, yj )Y |2 < ∞ . μ2j

and

(A.46)

j∈J

In this case x =

 1 (y, yj )Y xj μj j∈J

is a solution of (A.45). We note that the solvability conditions (A.46) require a fast decay of the Fourier coefficients of y with respect to the orthonormal system {yj : j ∈ J} in order for the series ∞  1 2 2 |(y, yj )Y | μ j=1 j to converge. Of course, this condition is only necessary for the important case where there exist infinitely many singular values. As a simple example, we study the following integral operator. Example A.59 Let K : L2 (0, 1) −→ L2 (0, 1) be defined by

t x(s) ds ,

(Kx)(t) :=

t ∈ (0, 1) , x ∈ L2 (0, 1) .

0

Then ∗

1

(K y)(t) =

y(s) ds

and

1  s

∗

(K Kx)(t) =

t

x(τ ) dτ ds . t

0

The eigenvalue problem K ∗ Kx = λx is equivalent to

1  s λx(t) =

 x(τ ) dτ ds ,

t

0



t ∈ [0, 1] .

A.7 The Fr´echet Derivative

351

Differentiating twice, we observe that for λ = 0 this is equivalent to the eigenvalue problem λx + x = 0 in (0, 1) , Solving this yields  xj (t) =

x(1) = x (0) = 0 .

2j − 1 2 cos πt, t ∈ [0, 1] , π 2

and λj =

4 (2j − 1)2 π 2

for j ∈ N. The singular values μj and the ONS {yj : j ∈ N} are given by μj

2 , j ∈ N, (2j − 1)π

=

 yj (t)

=

and

2j − 1 2 sin πt , j ∈ N . π 2

The singular value decomposition of Theorem A.57 !makes it possible to define, for every continuous function f : 0, K2L(X,Y ) → R, the operator f (K ∗ K) from X into itself by  f (K ∗ K)x = f (μ2j ) (x, xj )X xj , x ∈ X . (A.47) j∈J

This operator is always well-defined, linear, and bounded. It is compact if, and only√ if, f (0) = 0 (see Problem A.5). The special cases f (t) = t and f (t) √ = t are of particular importance. From this definition we note that R K ∗ K = R(K ∗ ). K of the previous Example A.59 we note that x ∈   For the operator R (K ∗ K)σ/2 if, and only if, ∞  j=1

1 |cj |2 < ∞ (2j − 1)2σ

 where

cj =

2 π

1 x(t) cos 0

2j − 1 πt dt 2

  are the Fourier coefficients of x. Therefore, R (K ∗ K)σ/2 plays the role of the σ periodic Sobolev spaces Hper (0, 2π) of Section A.4.

A.7

The Fr´ echet Derivative

In this section, we briefly recall some of the most important results for nonlinear mappings between normed spaces. The notions of continuity and differentiability carry over in a very natural way. Definition A.60 Let X and Y be normed spaces over the field K = R or C, U ⊂ X an open subset, x ˆ ∈ U , and T : X ⊃ U → Y be a (possibly nonlinear) mapping.

352

Appendix A

(a) T is called continuous at x ˆ if for every ε > 0 there exists δ > 0 such that T (x) − T (ˆ x)Y ≤ ε for all x ∈ U with x − x ˆX ≤ δ. (b) T is called Fr´echet differentiable at x ˆ ∈ U if there exists a linear bounded operator A : X → Y (depending on x ˆ) such that lim

h→0

1 T (ˆ x + h) − T (ˆ x) − AhY = 0 . hX

(A.48)

x) := A. In particular, T  (ˆ x) ∈ L(X, Y ). We write T  (ˆ (c) The mapping T is called continuously Fr´echet differentiable at x ˆ ∈ U if T is Fr´echet differentiable in a neighborhood V of x ˆ and the mapping T  : V → L(X, Y ) is continuous in x ˆ. Continuity and differentiability of a mapping depend on the norms in X and Y , in contrast to the finite-dimensional case. If T is differentiable in x ˆ, then the linear bounded mapping A in part (b) of Definition A.60 is unique. x) := A is well-defined. If T is differentiable in x, then T is also Therefore, T  (ˆ continuous in x. In the finite-dimensional case X = Kn and Y = Km , the linear bounded mapping T  (x) is given by the Jacobian (with respect to the Cartesian coordinates). Example A.61 (Integral Operator) Let f : [c, d] × [a, b] × K → K, f = f (t, s, r), K = R or K = C, be continuous with respect to all arguments and also ∂f /∂r ∈ C [c, d] × [a, b] × R . (a) Let the mapping T : C[a, b] → C[c, d] be defined by

b T (x)(t) :=

  f t, s, x(s) ds ,

t ∈ [c, d] , x ∈ C[a, b] .

(A.49)

a

We equip the normed spaces C[c, d] and C[a, b] with the maximum norm. Then T is continuously Fr´echet differentiable with derivative 





b

T (x)z (t) =

 ∂f  t, s, x(s) z(s) ds , ∂r

t ∈ [c, d] , x, z ∈ C[a, b] .

a

Indeed, let (Az)(t) be the term on the right-hand side. By assumption, ∂f /∂r is continuous on [a, b] × [c, d] × R, thus uniformly continuous on the compact set M = {(t, s,r) ∈ [a, b] × [c, d] × K : |r| ≤ x∞ + 1}. Let ε > 0. Choose   ∂f ε δ ∈ (0, 1) with  ∂f ˜) ≤ b−a for all (t, s, r), (t, s, r˜) ∈ M with ∂r (t, s, r) − ∂r (t, s, r

A.7 The Fr´echet Derivative

353

|r − r˜| ≤ δ. We estimate for z ∈ C[a, b] with z∞ ≤ δ:   T (x + z)(t) − T (x)(t) − (Az)(t)  b   $ %         ∂f  =  t, s, x(s) z(s) ds f t, s, x(s) + z(s) − f t, s, x(s) − ∂r   a  b 1         ∂f   d  =  f t, s, x(s) + rz(s) − t, s, x(s) z(s) dr ds dr ∂r   0

a

≤



b 1   ∂f   ∂f    |z(s)| dr ds  t, s, x(s) + rz(s) − t, s, x(s)   ∂r ∂r a

0

b ≤

z∞ a

ε ds = εz∞ . b−a

This holds for all t ∈ [c, d], thus also for the maximum. Therefore, T (x + z) − T (x) − Az∞ ≤ ε z∞

for all z ∈ C[a, b] with z∞ ≤ δ .

This proves the differentiability of T from C[a, b] into C[c, d].   (b) Let now in addition ∂f /∂r ∈ C [a, b] × [c, d] × K be Lipschitz continuous  with respect to r; that is, there exists κ > 0 with ∂f (t, s, r)/∂r − ∂f (t, s, r˜ ≤ κ|r − r˜| for all (t, s, r), (t, s, r˜) ∈ [a, b] × [c, d] × K. Then the operator T of (A.49) is also Fr´echet-differentiable as an operator from L2 (a, b) into C[c, d] (and thus into L2 (c, d)) with the same representation of the derivative. Indeed, define again the operator A as in part (a). Then A is bounded because

|(Az)(t)|

≤



b   ∂f    |z(s)| ds  t, s, x(s)   ∂r

≤

 %

b $  ∂f     + κ|x(s)| |z(s)| ds t, s, 0  ∂r 

a

a

≤

! max ∂f (τ, ·, 0)/∂rL2 (a,b) + κxL2 (a,b) zL2 (a,b)

c≤τ ≤d

for all t ∈ [c, d]. Concerning the computation of the derivative we can proceed in the same way

354

Appendix A

as in part (a):   T (x + z)(t) − T (x)(t) − (Az)(t) ≤



b 1   ∂f   ∂f    |z(s)| dr ds  t, s, x(s) + rz(s) − t, s, x(s)   ∂r ∂r 0

a

b 1 ≤

κ a

0

κ r |z(s)| dr ds = 2 2

b

|z(s)|2 ds =

κ z2L2 (a,b) . 2

a

This holds for all t ∈ [c, d]. Therefore, T (x + z) − T (x) − Az∞ ≤ which proves the differentiability.

κ 2

z2L2 (a,b)

The following theorem collects further properties of the Fr´echet derivative. Theorem A.62 (a) Let T, S : X ⊃ U → Y be Fr´echet differentiable at x ∈ U . Then T + S and λT are also Fr´echet differentiable for all λ ∈ K and (T + S) (x) = T  (x) + S  (x) ,

(λT ) (x) = λ T  (x) .

(b) Chain rule: Let T : X ⊃ U → V ⊂ Y and S : Y ⊃ V → Z be Fr´echet differentiable at x ∈ U and T (x) ∈ V , respectively. Then S T is also Fr´echet differentiable at x and   (S T ) (x) = S  T (x) T  (x) . /0 1 . /0 1

∈ L(X, Z) .

∈L(Y,Z) ∈L(X,Y )

(c) Special case: If x ˆ, h ∈ X are fixed and T : X → Y is Fr´echet differentiable on X, then ψ : K → Y , defined by ψ(t) := T (ˆ x + th), t ∈ K, x + th)h ∈ Y . Note that origiis differentiable on K and ψ  (t) = T  (ˆ nally ψ  (t) ∈ L(K, Y ). In this case, one identifies the linear mapping ψ  (t) : K → Y with its generating element ψ  (t) ∈ Y . If T  is Lipschitz continuous then the following estimate holds. Lemma A.63 Let T be differentiable in the ball B(x, ρ) centered at x ∈ U with radius ρ > 0, and let there exists γ > 0 with T  (x) − T  (x)L(X,Y ) ≤ γx − xX for all x ∈ B(x, ρ). Then   T (x) − T (x) − T  (x)(x − x) ≤ γ x − x2X Y 2

for all x ∈ B(x, ρ) .

A.7 The Fr´echet Derivative

355

Proof: Let  ∈ Y ∗ andx ∈ B(x, ρ) kept fixed. Set h = x − x and define the scalar function f (t) = , T (x + th) for |t| < ρ/hX where ·, · = ·, · Y ∗ ,Y denotes the dual pairing in Y ∗ ,Y . Then, by the chain rule of the previous    Theorem, f (t) = , T (x + th)h and thus     , T (x) − T (x) − T  (x)h    = , T (x) − , T (x) − , T  (x)h   1      = f (1) − f (0) − , T  (x)h  =  [f  (t) − , T  (x)h ] dt 0  1   !     =  , T (x + th)h − , T (x)h dt 0  1    !      =  , T (x + th) − T (x) h dt 0



≤

Y ∗ hX

≤

γY ∗ h2X

 T  (x + th) − T  (x) dt L(X,Y )

1

0



1

t dt = 0

γ Y ∗ h2X . 2

We set y := T (x) − T (x) − T  (x)h and choose  ∈ Y ∗ with Y ∗ = 1 and , y = yY . This is possible by a well known consequence of the Hahn-Banach theorem (see [151], Chap.V, §7, Theorem 2). Then the assertion follows.  We recall Banach’s contraction mapping principle (compare with Theorem A.31 for the linear case). Theorem A.64 (Contraction Mapping Principle) Let C ⊂ X be a closed subset of the Banach space X and T : X ⊃ C → X a (nonlinear) mapping with the properties (a) T maps C into itself; that is, T (x) ∈ C for all x ∈ C, and (b) T is a contraction on C; that is, there exists c < 1 with T (x) − T (y)X ≤ c x − yX

for all x, y ∈ C .

(A.50)

Then there exists a unique x ˜ ∈ C with T (˜ x) = x ˜. The sequence (x ) in C, ˜ for every x0 ∈ C. Furtherdefined by x+1 := T (x ),  = 0, 1, . . . converges to x more, the following error estimates hold: x+1 − x ˜X ≤ c x − x ˜X ,

 = 0, 1, . . . ;

(A.51a)

that is, the sequence converges linearly to x ˜, x − x ˜X

≤

x − x ˜X

≤

for  = 1, 2, . . .

c x1 − x0 X , (a priori estimate) 1−c 1 x+1 − x X , (a posteriori estimate) 1−c

(A.51b) (A.51c)

356

Appendix A

The Newton method for systems of nonlinear equations has a direct analogy for equations of the form T (x) = y, where T : X → Y is a continuously Fr´echet differentiable mapping between Banach spaces X and Y . We formulate a simplified Newton method and prove local linear convergence. It differs from the ordinary Newton method not only by replacing the derivative T  (x ) by x) but also by requiring only the existence of a left inverse. T  (ˆ Theorem A.65 (Simplified Newton Method) Let T : X → Y be continuously Fr´echet differentiable between Banach spaces X and Y . Let V ⊂ X be a closed subspace, x ˆ ∈ V and yˆ := T (ˆ x) ∈ Y . Let x) : X → Y L : Y → V be linear and bounded such that L is a left inverse of T  (ˆ x)v = v for all v ∈ V . on V ; that is, L T  (ˆ Then there exists ε > 0 such that for any y¯ = T (¯ x) with x ¯ ∈ X and ¯ x− ˜∈V: x ˆX ≤ ε the following algorithm converges linearly to some x x0 = x ˆ,

x+1 = x − L [T (x ) − y¯] ,

 = 0, 1, 2, . . . .

(A.52)

The limit x ˜ ∈ V satisfies L[T (˜ x) − y¯] = 0. Proof: We apply the contraction mapping principle of the preceding theorem to the mapping ! S(x) := x − L[T (x) − y¯] = L T  (ˆ x)x − T (x) + T (¯ x) from V into itself on some closed ball B[ˆ x, ρ] ⊂ V . We estimate S(x) − S(z)X

≤ ≤

LL(Y,X) T  (ˆ x)(x − z) + T (z) − T (x)Y 2 LL(Y,X) x − zX T  (ˆ x) − T  (z)L(X,Y ) T (z) − T (x) + T  (z)(x − z)Y 3 + x − zX

and S(x) − x ˆX

≤ ≤

LL(Y,X) T  (ˆ x)(x − x ˆ) − T (x) + T (¯ x)Y

LL(Y,X) T  (ˆ x)(x − x ˆ) + T (ˆ x) − T (x)Y + LL(Y,X) T (ˆ x) − T (¯ x)Y

First, we choose ρ > 0 such that $ % T (z) − T (x) + T  (z)(x − z)Y 1 LL(Y,X) T  (ˆ x) − T  (z)L(X,Y ) + ≤ x − zX 2 for all x, z ∈ B[ˆ x, ρ]. This is possible because T is continuously differentiable. Next, we choose ε > 0 such that LL(Y,X) T (ˆ x) − T (¯ x)Y ≤

ρ 2

A.8 Convex Analysis

357

for ¯ x−x ˆX ≤ ε. Then we conclude that S(x) − S(z)X

≤

S(x) − x ˆX

≤

1 x − zX for all x, z ∈ B[ˆ x, ρ] , 2 1 1 x − x ˆX + ρ ≤ ρ for all x ∈ B[ˆ x, ρ] . 2 2

Application of the contraction mapping principle ends the proof.



The notion of partial derivatives of mappings T : X × Z → Y is introduced just as for functions of two scalar variables as the Fr´echet derivative of the mappings T (·, z) : X → Y for z ∈ Z and T (x, ·) : Z → Y for x ∈ X. We denote the partial derivatives in (x, z) ∈ X × Z by ∂ T (x, z) ∈ L(X, Y ) and ∂x

∂ T (x, z) ∈ L(Z, Y ) . ∂z

Theorem A.66 (Implicit Function Theorem) Let T : X ×Z → Y be continuously Fr´echet differentiable with partial derivatives ∂ ∂ x, zˆ) = 0 ∂x T (x, z) ∈ L(X, Y ) and ∂z T (x, z) ∈ L(Z, Y ). Furthermore, let T (ˆ ∂ x, zˆ) : Z → Y be a norm-isomorphism from Z onto Y . Then there and ∂z T (ˆ exists a neighborhood U of x ˆ and a Fr´echet differentiable function ψ : U → Z such that ψ(ˆ x) = zˆ and T x, ψ(x) = 0 for all x ∈ U . The Fr´echet derivative ψ  ∈ L(X, Z) is given by $ %  −1 ∂   ∂  T x, ψ(x) T x, ψ(x) , x ∈ U . ψ  (x) = − ∂z ∂x The following special case is particularly important. ˆ = 0 and ∂ T (ˆ ˆ = 0. Let Z = Y = K; thus T : X × K → K and T (ˆ x, λ) x, λ) ∂λ Then there exists a neighborhood U of ˆ and a Fr´echet differentiable function  x ˆ and T x, ψ(x) = 0 for all x ∈ U and ψ : U → K such that ψ(ˆ x) = λ ψ  (x) = −

∂ ∂λ T

 ∂  1   T x, ψ(x) ∈ L(X, K) = X ∗ , x, ψ(x) ∂x

x∈U,

where again X ∗ denotes the dual space of X.

A.8

Convex Analysis

Definition A.67 Let X be a normed space. (a) A set M ⊂ X is called convex if for all x, y ∈ M and all λ ∈ [0, 1] also λx + (1 − λ)y ∈ M . (b) Let M ⊂ X be convex. A function f : M → R is called convex if for all x, y ∈ M and all λ ∈ [0, 1]   f λx + (1 − λ)y ≤ λf (x) + (1 − λ) f (y) .

358

Appendix A

(c) A function f : M → R is called concave if −f is convex; that is, if for all x, y ∈ M and all λ ∈ [0, 1]   f λx + (1 − λ)y ≥ λf (x) + (1 − λ) f (y) . (d) f is called strictly convex if   f λx + (1 − λ)y < λf (x) + (1 − λ) f (y) x, y ∈ M and all λ ∈ (0, 1) with x = y. The definition for a strictly concave function is formulated analogously. The definition of convexity of a set or a function can be extended to more than two elements. Lemma A.68 Let X be a normed space. (a) A set M ⊂ X is convex if, and m only if, for any elements xj ∈ M , j = 1, . . . , m, and λj ≥ 0 with j=1 λj = 1 also the convex combination m j=1 λj xj belongs to M . (b) Let M ⊂ X be convex. A function f : M → R is convex if, and only if, f

 m

 λj xj

≤

j=1

m 

λj f (xj )

j=1

m for all xj ∈ M , j = 1, . . . , m, and λj ≥ 0 with j=1 λj = 1. For concave functions the characterization holds analogously. For any set A ⊂ X of a normed space A the set conv A =

 m j=1

λj aj : aj ∈ A, λj ≥ 0,

m 

 λj = 1, m ∈ N

(A.53)

j=1

is convex by the previous lemma and is called the convex hull of A The following separation theorem is one of the most important tools in the area of convex analysis. Theorem A.69 Let X be a normed space over R and A, B ⊂ X two convex sets with A ∩ B = ∅. Furthermore, let A be open. Then there exists a hyperplane which separates A and B; that is, there exists  ∈ X ∗ and γ ∈ R such that  = 0 and , a X ∗ ,X ≥ γ ≥ , b X ∗ ,X for all a ∈ A and b ∈ B . Here, , x X ∗ ,X denotes the dual pairing; that is, the application of  ∈ X ∗ to x ∈ X.

A.8 Convex Analysis

359

For a proof we refer to,

e.g., [139], Chapter II. The hyperplane is given by x ∈ X : , x X ∗ ,X = γ . The convexity can be characterized easily for differentiable functions. Lemma A.70 Let X be a normed space and M ⊂ X be an open convex set and f : M → R be Fr´echet differentiable on M . Then f is convex if, and only if, f (y) − f (x) − f  (x)(y − x) ≥ 0

for all x, y ∈ M .

f is strictly convex if, and only if, the inequality holds strictly for all x = y. For concave functions the characterizations hold analogously. Proof: Let first f be  convex, x, y ∈ M and t ∈ (0, 1]. From ! the convexity of f we conclude that f x + t(y − x) ≤ f (x) + t f (y) − f (y) , thus f (y) − f (x) ≥ =

 ! 1  f x + t(y − x) − f (x) t  ! 1  f x + t(y − x) − f (x) − t f  (x)(y − x) + f  (x)(y − x) . t

The first term on the right-hand side tends to zero as t → 0. This proves f (y) − f (x) ≥ f  (x)(y − x). Let now f (u) − f (v) ≥ f  (v)(u − v) for all u, v ∈ M . For x, y ∈ M and λ ∈ [0, 1] apply this twice to v = λx + (1 − λ)y and u = y and u = x, respectively. With y − v = −λ(x − y) and x − v = (1 − λ)(x − y) this yields f (y) − f (v) ≥

f  (v)(y − v) = λf  (v)(y − x) ,

f (x) − f (v) ≥

f  (v)(x − v) = (1 − λ)f  (v)(x − y) .

Multiplying the first inequality by 1 − λ and the second by λ and adding these inequalities yields the assertion.  Note that we could equally well write f  (x), y − x X ∗ ,X for f  (x)(y − x). We use both notations synonymously. This characterization motivates the definition of the subgradient of a convex function. Definition A.71 Let X be a normed space over R with dual X ∗ , M ⊂ X be an open convex set, and f : M → R be a convex function. For x ∈ M the set

∂f (x) :=  ∈ X ∗ : f (z) − f (x) − , z − x X ∗ ,X ≥ 0 for all z ∈ M is called the subgradient of f at x. As one sees from the function f (x) = |x| for x ∈ R the subgradient ∂f is, in general, a multivalued function. It is not empty for continuous functions, and it consists of the derivative as the only element for differentiable functions.

360

Appendix A

Lemma A.72 Let X be a normed space over R with dual X ∗ , M ⊂ X be an open convex set, and f : M → R be a convex and continuous function. (a) Then ∂f (x) = ∅ for all x ∈ M .

(b) If f is Fr´echet differentiable at x ∈ M then ∂f (x) = f  (x) . In particular, ∂f is single valued at x.

Proof: (a) Define the set D ⊂ X × R by

D := (z, r) ∈ M × R : r > f (z) . Then D is open because M is open and f is continuous and  D is also convex (see Problem A.6). Fix x ∈ M . Then we observe that x, f (x) ∈ / D. The separation theorem for convex sets (see Theorem A.69) yields the existence of (, s) ∈ X ∗ × R with (, s) = (0, 0) and γ ∈ R such that , z X ∗ ,X + sr ≤ γ ≤ , x X ∗ ,X + s f (x)

for all r > f (z) , z ∈ M .

Letting r tend to +∞ implies that s ≤ 0. Also s = 0 because otherwise , z X ∗ ,X ≤ , x X ∗ ,X for all z ∈ M which would imply that also  vanishes2 , a contradiction to (, s) = (0, 0). Therefore, s < 0 and, without loss of generality (division by |s|), we can assume that s = −1. Letting r tend to f (z) yields , z X ∗ ,X − f (z) ≤ , x X ∗ ,X − f (x) which shows that  ∈ ∂f (x). (b) f  (x) ∈ ∂f (x) follows from Lemma A.70. Let  ∈ ∂f (x) and y ∈ X arbitrary with y = 0. For sufficiently small t > 0 we have that x + ty ∈ M and thus   f (x + ty) − f (y) − tf  (x)y ≥ t  − f  (x) y . Division by t > 0 and letting t tend to zero implies thatthe left-hand  side tends to zero by the definition of the derivative. Therefore,  − f  (x) y ≤ 0. Since  this holds for all y ∈ X we conclude that  = f  (x). Lemma A.73 Let ψ : [0, a] → R be a continuous, concave, and monotonically increasing function with ψ(0) = 0. (a) Then ψ(st) ≤ max{1, s} ψ(t) for all s, t ≥ 0 with st, t ∈ [0, a]. √ !2 (b) The function t → ψ( t) is concave on [0, a2 ]. (c) Let K : X → Y be a linear compact operator between Hilbert spaces and a ≥ KL(X,Y ) . Then   ∗ 1/2     ψ (K K) z X ≤ ψ KzY

2 The

for all z ∈ X with zX ≤ 1 .

reader should verify this himself by using that M is an open set.

A.8 Convex Analysis

361

Proof: (a) If s ≤ 1 the assertion follows from the monotonicity of ψ. If s ≥ 1 then       1 1 1 1 1 (st) + 1 − ψ(st) + 1 − ψ(st) . ψ(t) = ψ 0 ≥ ψ(0) = s s s s s √ !2 (b) Set φ(t) = ψ( t) for t ∈ [0, a2 ]. If ψ was twice differentiable then an elementary calculation shows that √ √ √ ! √ ψ  ( t) √  √ 1   φ (t) = ψ( t ψ ( t) − ψ( t) + t) ψ ( t) . 2t 2t3/2 The first term is non-positive because ψ  ≥ 0 and s ψ  (s) − ψ(s) = ψ(0) − ψ(s) − (0 − s)ψ  (s) ≤ 0 by Lemma A.70. The second term is also non-positive because ψ ≥ 0 and ψ  ≤ 0. Therefore, φ (t) ≤ 0 for all t which proves that φ is concave. If ψ is merely continuous then we approximate ψ by a sequence (ψ ) of smooth concave and monotonically increasing functions with ψ (0) = 0 and ψ − ψ∞ → 0 as  tends to infinity. We sketch the proof but leave the details to the reader. In the first step we approximate ψ by the interpolating polygonal function pm a with respect to tj = j m , j = 0, . . . , m, with values ψj := ψ(tj ) at tj . Then, for any  ∈ N there exists m = m() ∈ N with pm − ψ∞ ≤ 1/. Next, we extend pm onto R by extending the first and the last segment linearly; that is, m−1 by setting pm (t) = ψ1 tt1 for t ≤ 0 and pm (t) = ψ(a) + ψ(a)−ψ a−tm−1 (t − a) for t ≥ a. In the third step we smooth pm by using a mollifier;that is, a non-negative 1 function φ ∈ C ∞ (R) with φ(t) = 0 for |t| ≥ 1 and −1 φ(t)dt = 1. We set & ' φρ (t) = ρ1 φ ρt and

∞ ψ˜ρ (t) =

ρ φρ (t − s) pm (s) ds =

−∞

φρ (s) pm (t − s) ds ,

t ∈ [0, a] .

−ρ

Then ψ˜ρ is in C ∞ [0, a], concave, monotonically increasing and ψ˜ρ −pm C[0,a] → 0 as ρ tends to zero. Finally we set ψ (t) = ψ˜ρ (t) − ψ˜ρ (0), where ρ = ρ() > 0 is such that ψ − pm C(0,a) ≤ 1/. This sketches the construction of the sequence ψ . √ !2 From the first part we know that t → ψ ( t) is concave for every . Letting  tend to infinity proves the assertion. √ !2 ! (c) We set again φ(t) = ψ( t) for t ∈ 0, K2L(X,Y ) and let z ∈ X with zX ≤ 1. We decompose z = z0 + z ⊥ with z0 ∈ N (K) and z ⊥ ⊥  ∗ 1/2 z in the form ∗ z = ψ (K K)1/2 z ⊥ and Kz = Kz ⊥ . Therefore, it N (K). Then ψ (K K) suffices to take z ∈ X with z ⊥ N (K) and zX ≤ 1. Witha singular system {μj , xj , yj : j ∈ J} of K we expand such a z ∈ X as z = j∈J zj xj . We set  Jn = J if J is finite and Jn = {1, . . . , n} if J = N and zˆ(n) = j∈Jn zˆj xj with

362

Appendix A

4  2 zˆj = zj / z (n) 2X = j∈Jn |ˆ zj |2 = 1 and thus ∈Jn |z | . Then ˆ   ∗ 1/2  (n) 2 ψ (K K) zˆ X

=



 !2 2 ψ(μj ) |ˆ zj | = φ(μ2j ) |ˆ zj | 2

j∈Jn

≤

φ



j∈Jn

μ2j |ˆ zj | 2



⎡ ⎛ = ⎣ψ ⎝

j∈Jn

"  # 2 = ψ K zˆ(n) Y



⎞⎤2 μ2j |ˆ zj | 2 ⎠ ⎦

j∈Jn

(n) where we used that  φ  is ∗concave.  Letting n tend to infinity yields zˆ → zˆ = 1/2   zˆ X ≤ ψ K zˆY . Therefore, z/zX and thus ψ (K K)   ∗ 1/2       ψ (K K) z X = zX ψ (K ∗ K)1/2 zˆX     1 ≤ zX ψ K zˆY = zX ψ KzY zX   ≤ ψ KzY

where we used the estimate ψ(st) ≤ s ψ(t) from part (a) for s = 1/zX ≥ 1  and t = KzY . Lemma A.74 There exist constants c+ > 0 and cp ≥ 0 with cp = 0 for p ≤ 2 and cp > 0 for p > 2 such that for all z ∈ R  c+ |z|p if p ≤ 2 or |z| ≥ 12 , p p cp |z| ≤ |1 + z| − 1 − pz ≤ (A.54) c+ |z|2 if p > 2 and |z| ≤ 12 . Proof: To show the upper estimate, let first |z| ≥ 12 . Then 1 ≤ 2|z| and thus   |1+z|p −1−pz ≤ (1+|z|)p +p·1·|z| ≤ (3|z|)p +p 2p−1 |z|p−1 |z| = 3p +p 2p−1 |z|p . p p Let now |z| ≤ 12 ; that is, 1 ≥ 2|z|. Then := f (z). ! |1+z| −1−pz = (1+z) −1−pz  p−1 − 1 and f (z) = p(p − 1)(1 + z)p−2 . Taylor’s We compute f (z) = p (1 + z) formula yields

(1 + z)p − 1 − pz = f (z) = f (0) + f  (0)z +

1  p(p − 1) f (ξ)z 2 = (1 + ξ)p−2 z 2 2 2

for some ξ with |ξ| ≤ |z|. Let first p ≥ 2. Then (1 + ξ)p−2 ≤ (3/2)p−2 , and the estimate is shown. 1 1 p−2 because 1 + ξ ≥ Let now p < 2. Then (1 + ξ)p−2 = (1+ξ) 2−p ≤ |z|2−p = |z| 1 − |ξ| ≥ 2|z| − |z| = |z|. Therefore, f (z) ≤ the upper estimate of (A.54) is shown.

p(p−1) 2

|z|p−2 z 2 =

p(p−1) 2

|z|p , and

(p−1)/(p−2)

To show the lower estimate, we choose cp ∈(0, 1) for p>2 such that cp −  1/(p−2) p−1 ≥ 0 and set cp = 0 for p ≤ 2. We distinguish between cp + 1 − cp

A.9 Weak Topologies

363

several cases. Case A: z ≥ 0. Set f (z) := |1 + z|p − 1 − pz − cp |z|p = (1 + z)!p − 1 − pz − cp z p for z ≥ 0. Then f (0) = 0 and f  (z) = p (1! + z)p−1 − 1 − cp z p−1 and !f  (0) = 0 and f  (z) = p(p − 1) (1 + z)p−2 − cp z p−2 ≥ p(p − 1) z p−2 − cp z p−2 ≥ 0 because cp ≤ 1. Therefore, f  is monotonically increasing, thus positive. Therefore, f is monotonically increasing, thus positive. Case B: z ≤ 0. We replace z by −z ≥ 0 and have to show that f (z) := |1 − z|p − 1 + pz − cp z p ≥ 0 for all z ≥ 0. p − 1 + pz − cp z p and f (0) = 0 and Case B1: 0 ≤ z ≤ 1. Then f (z) = ! (1 − z)  p−1 p−1  and f (0) = 0 and f  (1) = p(1 − cp ) > 0. − cp z f (z) = p 1 − (1 − z)  If p ≤ 2 then cp = 0 and thus f ≥ 0, thus f ≥ 0 on [0, 1].! Let p > 2. Then f  (z) = p(p − 1) (1 − z)p−2 − cp z p−2 and f  (0) > 0 and f  (1) < 0. Since f  < 0 on (0, 1) there exists exactly one zero zˆ ∈ (0, 1) of f  . z , 1]. Therefore, the minimal Therefore, f  increases on [0, zˆ] and decreases on [ˆ values of f  on [0, 1] are obtained for z = 0 or z = 1 which are both non-negative. Therefore, f  ≥ 0 on [0, 1] which implies that also f is non-negative.

−cp z p and f (1) = −1+p−cp > 0 Case B2: z ≥ 1. Then f (z) = (z −1)p −1+pz !  p−1 p−1 and f  (1) = p(1 − cp ) > 0. + 1 − cp z and f (z) = p (z − 1)  If p ≤ 2 we conclude that f ≥ 0 on [1, ∞), thus f  ≥ 0 on [1, ∞) and thus also f ≥ 0 on [1, ∞). ! Let finally p > 2. Then f  (z) = p(p − 1) (z − 1)p−2 − cp z p−2 and f  (1) < 0 and f  (z) → ∞ as z → ∞. The second derivative f  vanishes at zˆ = 1 − 1/(p−2) !−1 cp > 1, which is negative on [1, zˆ) and positive for z > zˆ. Therefore, f  (z) decreases on (1, zˆ) and increases for z > zˆ. Its minimal value is attained at zˆ which is computed as ! f  (ˆ z ) = p (ˆ z − 1)p−1 − cp zˆp−1 + 1 p−1 !  p c(p−1)/(p−2) ≥ 0 − cp + 1 − cp1/(p−2) =  p 1/(p−2) p−1 1 − cp by the choice of cp . Therefore, f  is positive for z ≥ 1 and thus also f .

A.9



Weak Topologies

In this subsection, we recall the basics on weak topologies. We avoid the definition of the topology itself; that is, the definition of the family of open sets, but restrict ourselves to the concept of weak convergence (or weak∗ convergence). Definition A.75 Let X be a normed space and X ∗ its dual space with dual pairing , x X ∗ ,X for  ∈ X ∗ and x ∈ X. (a) A sequence (xn) in X is said to converge weakly to x∈X if lim , xn X ∗ ,X = n→∞

, x X ∗ ,X for all  ∈ X ∗ . We write xn  x for the weak convergence.

364

Appendix A

(b) A sequence (n ) in X ∗ is said to converge weak∗ to  ∈ X ∗ if ∗ lim n , x X ∗ ,X = , x X ∗ ,X for all x ∈ X. We write n   for the n→∞ weak∗ convergence. First we note that weak convergence is indeed weaker than norm convergence. This follows directly from the continuity of the functional  ∈ X ∗ . Furthermore, we note that weak convergence in X ∗ means that T, n X ∗∗ ,X ∗ → T,  X ∗∗ ,X ∗ as n → ∞ for all T ∈ X ∗∗ where now ·, · X ∗∗ ,X ∗ denotes the dual pairing in (X ∗∗ , X ∗ ). By the canonical imbedding X → X ∗∗ (see formula (A.14)) weak convergence in X ∗ is stronger than weak∗ convergence. In reflexive spaces (see Definition A.22) the bidual X ∗∗ can be identified with X and, therefore, weak and weak∗ convergence coincide. If X is a Hilbert space with inner product (·, ·)X then, by the representation Theorem A.23 of Riesz, the dual space X ∗ is identified with X itself, and weak convergence xn  x is equivalent to (xn , z)X → (x, z)X for all z ∈ X. We will need the following results which we cite without proof. Theorem A.76 Let X be a normed space and X ∗ its dual space and (xn ) a sequence in X which converges weakly to some x ∈ X. Then the following holds: (a) The weak limit is well-defined; that is, if xn  x and xn  y then x = y. (b) The sequence (xn ) is bounded in norm; that is, xn X is bounded in R. (c) If X is finite dimensional then (xn ) converges to x in norm. Therefore, in finite-dimensional spaces there is no difference between weak and norm convergence. (d) Let U ⊂ X be a convex and closed set. Then U is also weakly sequentially closed; that is, every weak limit point x ∈ X of a weakly convergent sequence (xn ) in U also belongs to U . (e) Let Y be another normed space and K : X → Y be a linear bounded operator. Then (Kxn ) converges weakly to Kx in Y . If K is compact then (Kxn ) converges in norm to Kx. The following theorem of Alaoglu–Bourbaki (see, e.g., [139]) is the essential ingredient to assure the existence of global minima of the Tikhonov functional for nonlinear inverse problems. Theorem A.77 Let X be a Banach space with dual space X ∗ . Then every bounded sequence (n ) in X ∗ has a weak∗ convergent subsequence. For Hilbert spaces the theorem takes the following form. Corollary A.78 Let X be a Hilbert space. Every bounded sequence (xn ) in X contains a weak accumulation point; that is, a weakly convergent subsequence.

A.10 Problems

A.10

365

Problems

A.1 (a) Show with Definition A.5 that a set M ⊂ X is closed if, and only if, the limit of every convergent sequence (xk )k in M also belongs to M . (b) Show that a subspace V ⊂ X is dense if, and only if, the orthogonal complement is trivial; that is, V ⊥ = {0}. A.2 Try to prove Theorem A.8. A.3 Let V ⊂ X ⊂ V  be a Gelfand triple (see Definition A.26) and j : V → X and j  : X → V  the corresponding embedding operators. Show that both of them are one-to-one with dense range. A.4 Define the operator A from the sequence space 2 into itself by  0, if k = 1 , (Ax)k := xk−1 , if k ≥ 2 , for x = (xk ) ∈ 2 . Show that λ = 1 is not an eigenvalue of A but I − A fails to be surjective. Therefore, 1 belongs to the spectrum of A. A.5 Let K : X → Y a compact operator between Hilbert spaces X and Y , and ! let f : 0, K2L(X,Y ) → R be continuous. Define the operator f (K ∗ K) from X into itself by (A.47). Show that this operator is well defined (that is, f (K ∗ K)x defines an element in X for every x ∈ X), linear, and bounded and that it is compact if, and only if, f (0) = 0. A.6 Let A ⊂ X be an open set of a Hilbert space X and f : A → R convex and continuous. Show that the set D ⊂ X × R, defined by

D := (z, r) ∈ A × R : r > f (z) . is convex and open in X × R.

Appendix B

Proofs of the Results of Section 2.7 In this appendix, we give the complete proofs of the theorems and lemmas of Chapter 2, Section 2.7. For the convenience of the reader, we formulate the results again. Theorem 2.20 (Fletcher–Reeves) Let K : X → Y be a bounded, linear, and injective operator between Hilbert spaces X and Y . The conjugate gradient method is well-defined and either stops or produces sequences (xm ), (pm ) in X with the properties   ∇f (xm ), ∇f (xj ) X = 0 for all j = m , (2.41a) and (Kpm , Kpj )Y = 0

for all j = m ;

(2.41b)

that is, the gradients are orthogonal and the directions pm are K-conjugate. Furthermore,   (2.41c) ∇f (xj ), K ∗ Kpm X = 0 for all j < m . Proof:

First, we note the following identities:

(α) ∇f (xm+1 ) = 2 K ∗ (Kxm+1 − y) = 2 K ∗ (Kxm − y) − 2tm K ∗ Kpm = ∇f (xm ) − 2tm K ∗ Kpm .     (β) pm , ∇f (xm+1 ) X = pm , ∇f (xm ) X − 2tm (Kpm , Kpm )Y = 0 by the definition of tm .   1 m 2 m 2 (γ) tm = 12 ∇f (xm ), pm X Kpm −2 Y = 4 ∇f (x )X /Kp Y since pm = 12 ∇f (xm ) + γm−1 pm−1 and (β). Now we prove the following identities by induction with respect to m:

© Springer Nature Switzerland AG 2021 A. Kirsch, An Introduction to the Mathematical Theory of Inverse Problems, Applied Mathematical Sciences 120, https://doi.org/10.1007/978-3-030-63343-1

367

368

Appendix B

  (i) ∇f (xm ), ∇f (xj ) X = 0 for j = 0, . . . , m − 1, (ii) (Kpm , Kpj )Y = 0 for j = 0, . . . , m − 1. Let m = 1. Then, using (α),     (i) ∇f (x1 ), ∇f (x0 ) X = ∇f (x0 )2X − 2t0 Kp0 , K∇f (x0 ) Y = 0 , which vanishes by (γ) since p0 = 12 ∇f (x0 ). (ii) By the definition of p1 and identity (α), we conclude that (Kp1 , Kp0 )Y

=

(p1 , K ∗ Kp0 )X

1 = − 2t0 1 = − 2t0

$

%  1 1 0 1 0 ∇f (x ) + γ0 p , ∇f (x ) − ∇f (x ) X 2   ! 1 ∇f (x1 )2X − γ0 p0 , ∇f (x0 ) X = 0 , 2

where we have used (β), the definition of p0 , and the choice of γ0 . Now we assume the validity of (i) and (ii) for m and show it for m + 1: (i) For j = 0, . . . m − 1 we conclude that (setting γ−1 = 0 in the case j = 0)     ∇f (xm+1 ), ∇f (xj ) X = ∇f (xm ) − 2tm K ∗ Kpm , ∇f (xj ) X   = −2tm ∇f (xj ), K ∗ Kpm X   = −4tm Kpj − γj−1 Kpj−1 , Kpm Y = 0 , where we have used 12 ∇f (xj ) + γj−1 pj−1 = pj and assertion (ii) for m. For j = m, we conclude that   ∇f (xm+1 ), ∇f (xm ) X   = ∇f (xm )2X − 2tm ∇f (xm ), K ∗ Kpm X  1 ∇f (xm )2X  ∇f (xm ), K ∗ Kpm X = ∇f (xm )2X − 2 m 2 Kp Y by (γ). Now we write





 1 m m−1 ∇f (x ) + γm−1 p (Kp , Kp )Y = Kp , K 2 Y  1 m Kp , K∇f (xm ) Y , = 2   which implies that ∇f (xm+1 ), ∇f (xm ) X vanishes. (ii) For j = 0, . . . , m − 1, we conclude that, using (α),   1 ∇f (xm+1 ) + γm pm , K ∗ Kpj (Kpm+1 , Kpj )Y = 2 X  1  m+1 j+1 ∇f (x = − ), ∇f (x ) − ∇f (xj ) X , 4tj m

m

m

Appendix B

369

which vanishes by (i). For j = m by (α) and the definition of pm+1 , we have   1 1 (Kpm+1 , Kpm )Y = ∇f (xm+1 ) + γm pm , ∇f (xm ) − ∇f (xm+1 ) 2tm 2 X   1 1 1 ∇f (xm+1 ), ∇f (xm ) X − ∇f (xm+1 )2X = 2tm 2 2   m  m   m m+1 + γm p , ∇f (x ) X − γm p , ∇f (x ) X . /0 1 = 12 ∇f (xm ) 2X

1 γm ∇f (xm )2X − ∇f (xm+1 )2X 4tm

=

by (i) and (β). This term vanishes by the definition of γm . Thus we have proven (i) and (ii) for m + 1 and thus for all m = 1, 2, 3 . . . To prove (2.41c) we write     ∇f (xj ), K ∗ Kpm X = 2 pj − γj−1 pj−1 , K ∗ Kpm X = 0 for j < m , and note that we have already shown this in the proof.



Theorem 2.21 Let (xm ) and (pm ) be the sequences of the conjugate gradient method. Define the space Vm := span {p0 , . . . , pm }. Then we have the following equivalent characterizations of Vm :

(2.42a) Vm = span ∇f (x0 ), . . . , ∇f (xm ) =



span p0 , K ∗ Kp0 , . . . , (K ∗ K)m p0

(2.42b)

for m = 0, 1, . . . . Furthermore, xm is the minimum of f (x) = Kx − y2Y on Vm−1 for every m ≥ 1.

Proof: Let V˜m = span ∇f (x0 ), . . . , ∇f (xm ) . Then V0 = V˜0 . Assume that we have already shown that Vm = V˜m . Since pm+1 = 12 ∇f (xm+1 ) + γm pm we also have that Vm+1 = V˜m+1 . Analogously, define the space Vˆm := span p0 , . . . , (K ∗ K)m p0 . Then V0 = Vˆ0 and V1 = Vˆ1 . Assume that we have already shown that Vj = Vˆj for all j = 0, . . . , m. Then we conclude that pm+1

= K ∗ (Kxm+1 − y) + γm pm = K ∗ (Kxm − y) − tm K ∗ Kpm + γm pm = pm − γm−1 pm−1 − tm K ∗ Kpm + γm pm ∈ Vˆm+1 .

(B1)

On the other hand, from (K ∗ K)m+1 p0 = (K ∗ K) (K ∗ K)m p0

!

∈ (K ∗ K)(Vm )

and K ∗ Kpj ∈ Vj+1 ⊂ Vm+1 by (B1) for j = 0, . . . , m, we conclude also that Vm+1 = Vˆm+1 .

370

Appendix B

Now every xm lies in Vm−1 . This is certainly true for m = 1, and if it holds for m then it holds also for m + 1 since xm+1 = xm − tm pm ∈ Vm . xm is the minimum of f on Vm−1 if and only if (Kxm − y, Kz)Y = 0 for all z ∈ Vm−1 . By (2.42a), this is the case if and only if ∇f (xm ), ∇f (xj ) X = 0 for all j = 0, . . . , m − 1. This holds by the preceding theorem.  Lemma 2.22 (a) The polynomial Qm , defined by Qm (t) = 1 − t Pm−1 (t) with Pm−1 from (2.43), minimizes the functional H(Q) := Q(KK ∗ )y2Y and satisfies

on

{Q ∈ Pm : Q(0) = 1}

  H Qm = Kxm − y2Y .

(b) For k = , the following orthogonality relation holds: Qk , Q :=

∞ 

 2 μ2j Qk (μ2j ) Q (μ2j ) (y, yj )Y  = 0 .

(2.44)

j=1

If y ∈ / span {y1 , . . . , yN } for any N ∈ N, then ·, · defines an inner product on the space P of all polynomials. Proof:  (a) Let Q ∈ Pm be an arbitrary polynomial with Q(0) = 1. Set P(t) := 1 − Q(t) /t and x := P(K ∗ K)K ∗ y = −P(K ∗ K)p0 ∈ Vm−1 . Then y − Kx = y − KP(K ∗ K)K ∗ y = Q(KK ∗ )y . Thus

    H Q = Kx − y2Y ≥ Kxm − y2Y = H Qm .

(b) Let k = . From the identity ∞

 1 ∇f (xk ) = K ∗ (Kxk − y) = − μj Qk (μ2j )(y, yj )Y xj , 2 j=1 we conclude that 0 =

 1 ∇f (xk ), ∇f (x ) X = Qk , Q . 4

The properties of the inner product are obvious, except perhaps the definiteness. If Qk , Qk = 0, then Qk (μ2j ) (y, yj )Y vanishes for all j ∈ N. The assumption on y implies that (y, yj )Y = 0 for infinitely many j. But then the polynomial Qk has infinitely many zeros μ2j . This implies Qk = 0, which ends the proof.  The following lemma is needed for the proof of Theorem 2.24. Lemma Let 0 < m ≤ m(δ) where m(δ) satisfies (2.46) and define the space X σ = (K ∗ K)σ/2 for any σ > 0, equipped with the norm xX σ :=(K ∗ K)−σ/2 xX .

Appendix B

371

Let y δ ∈ / span{y1 , . . . , yN } for any N ∈ N with y δ − Kx∗ Y ≤ δ, and let σ x ∈ X for some σ > 0, and x∗ X σ ≤ E. Then ∗

E Kxm,δ − y δ Y ≤ δ + (1 + σ)(σ+1)/2    d Qδm (0)(σ+1)/2 dt

where Qδm denotes the polynomial Qm for y = y δ . Before we prove this lemma, we recall some properties of orthogonal functions (see [254]). As we saw in Lemma 2.22, the polynomials Qm are orthogonal with respect to the inner product ·, · . Therefore, the zeros λj,m , j = 1, . . . , m, of Qm are all real and positive and lie in the interval 0, K2L(X,Y ) . By their normalization, Qm must have the form m  9

Qm (t) =

j=1

1−

t λj,m

 .

Furthermore, the zeros of two subsequent polynomials interlace; that is, 0 < λ1,m < λ1,m−1 < λ2,m < λ2,m−1 < · · · < λm−1,m−1 < λm,m < K2L(X,Y ) . Finally, from the factorization of Qm , we see that  d 1 Qm (t) = −Qm (t) and dt λ −t j=1 j,m $ % 2  m m d2 1 1 Q (t) = Q (t) − m m dt2 λ −t (λj,m − t)2 j=1 j,m j=1 m

= Qm (t)

m  j,=1 j=

1 . (λj,m − t)(λ,m − t) 2

d d For 0 ≤ t ≤ λ1,m , we conclude that dt Qm (t) ≤ 0, dt 2 Qm (t) ≥ 0, and 0 ≤ Qm (t) ≤ 1. For the proof of the lemma and the following theorem, it is convenient to introduce two orthogonal projections. For any ε > 0, we denote by Lε : X → X and Mε : Y → Y the orthogonal projections  Lε z := (z, xn )X xn , z ∈ X , μ2n ≤ε

Mε z

:=



(z, yn )Y yn ,

z∈Y

μ2n ≤ε

where {μn , xn , yn : n ∈ J} denotes a singular system of K.

372

Appendix B The following estimates are easily checked: √

Mε KxY ≤

1 and (I − Lε )xX ≤ √ (I − Mε )KxY ε

εLε xX

for all x ∈ X. Proof of the lemma: Let λj,m be the zeros of Qδm . We suppress the dependence on δ. The orthogonality relation (2.44) implies that Qδm is orthogonal to the polynomial t → Qδm (t)/(λ1,m − t) of degree m − 1; that is, ∞ 

μ2n Qδm (μ2n )

n=1

2 Qδm (μ2n )  δ (y , yn )Y  = 0 . 2 λ1,m − μn

This implies that 

Qδm (μ2n )2

μ2n ≤λ1,m

= ≥



 δ  μ2n (y , yn )Y 2 2 λ1,m − μn

 δ  μ2n (y , yn )Y 2 μ2n − λ1,m μ2n >λ1,m   2 Qδm (μ2n )2 (y δ , yn )Y  . Qδm (μ2n )2

μ2n >λ1,m

From this, we see that ⎛ y δ − Kxm,δ 2Y



= ⎝

+

μ2n ≤λ1,m



≤

μ2n ≤λ1,m

⎞



  ⎠ Qδm (μ2n )2 (y δ , yn )Y 2

μ2n >λ1,m



Qδm (μ2n )2 .

  δ  μ2n (y , yn )Y 2 1+ 2 λ1,m − μn /0 1

=:Φm (μ2n )2

= Mλ1,m Φm (KK ∗ )y δ 2Y , where we have set 

 Φm (t) :=

Qδm (t)

1+

t λ1,m − t

=

Qδm (t)

λ1,m . λ1,m − t

Therefore, y δ − Kxm,δ Y ≤ Mλ1,m Φm (KK ∗ )(y δ − y ∗ )Y + Mλ1,m Φm (KK ∗ )Kx∗ Y .

Appendix B

373

We estimate both terms on the right-hand side separately:   2 Mλ1,m Φm (KK ∗ )(y δ − y ∗ )2Y = Φm (μ2n )2 (y δ − y ∗ , yn )Y  μ2n ≤λ1,m

≤ Mλ1,m Φm (KK ∗ )Kx∗ 2Y

=

max

Φm (t)2 y δ − y ∗ 2Y ,



2 ! −2σ  ∗ Φm (μ2n )2 μ2+2σ μn (x , xn )X  n

0≤t≤λ1,m

μ2n ≤λ1,m

≤

max

0≤t≤λ1,m

! t1+σ Φm (t)2 x∗ 2X σ .

The proof is finished provided we can show that 0 ≤ Φm (t) ≤ 1 and σ+1  1+σ 1+σ 2 t Φm (t) ≤ for all 0 ≤ t ≤ λ1,m . d δ | dt Qm (0)| The first assertion follows from Φm (0) = 1, Φm (λ1,m ) = 0, and ! d Φm (t)2 = dt

λ1,m d λ1,m + Qm (t)2 2 Qm (t) Qm (t) dt λ1,m − t (λ1,m − t)2 ⎤ ⎡ m  1 1 2⎣ ⎦ ≤ 0. −2 Φm (t) λ1,m − t λ − t j,m j=1

=

Now we set ψ(t) := t1+σ Φm (t)2 . Then ψ(0) = ψ(λ1,m ) = 0. Let tˆ ∈ (0, λ1,m ) be the maximum of ψ in this interval. Then ψ  (tˆ) = 0, and thus by differentiation ! d (σ + 1) tˆσ Φm (tˆ)2 + tˆσ+1 Φm (tˆ)2 = 0 ; dt that is,

⎡

≥

⎤

m  1 ⎦ ≥ tˆ ˆ ˆ λ −t λ1,m − t λ − tˆ j=1 j,m j=1 j,m   m  d δ  1  ˆ ˆ = t  Qm (0) . t λ dt j=1 j,m

σ + 1 = tˆ⎣2

m 

1

−

1

d δ This implies that tˆ ≤ (σ + 1)/| dt Qm (0)|. With ψ(t) ≤ tˆσ+1 for all t ∈ [0, λ1,m ], the assertion follows. 

Theorem 2.24 Assume that y ∗ and y δ do not belong to the linear span of finitely many yj . Let the sequence xm(δ),δ be constructed by the conjugate gradient method with stopping rule (2.46) for fixed parameter r > 1. Let x∗ = (K ∗ K)σ/2 z ∈ X σ for some σ > 0 and z ∈ X. Then there exists c > 0 with x∗ − xm(δ),δ X ≤ c δ σ/(σ+1) E 1/(σ+1) ,

(2.47)

374

Appendix B

where E = zX . Proof: Similar to the analysis of Landweber’s method, we estimate the error by the sum of two terms: the first converges to zero as δ → 0 independently of m, and the second term tends to infinity as m → ∞. The role of the norm d δ 1/2 Qm (0) . Rα L(Y,X) here is played by  dt First, let δ and m := m(δ) be fixed. Set for abbreviation d  q :=  Qδm (0) . dt Choose 0 < ε ≤ 1/q ≤ λ1,m . With x ˜ := x∗ − Qδm (K ∗ K)x∗ = Pδm−1 (K ∗ K)K ∗ y ∗ we conclude that x∗ − xm,δ X

≤ ≤

Lε (x∗ − xm,δ )X + (I − Lε )(x∗ − xm,δ )X Lε (x∗ − x ˜X + Lε (˜ x − xm,δ )X 1 + √ (I − Mε )(y ∗ − Kxm,δ )Y ε

≤

Lε Qδm (K ∗ K)x∗ X + Lε Pδm−1 (K ∗ K)K ∗ (y ∗ − y δ )X 1 1 + √ y ∗ − y δ Y + √ y δ − Kxm,δ Y ε ε    σ/2 δ √ 1+r E max t Qm (t) + δ max  t Pδm−1 (t) + √ δ . 0≤t≤ε 0≤t≤ε ε

≤

From ε ≤ λ1,m and 0 ≤ Qδm (t) ≤ 1 for 0 ≤ t ≤ λ1,m , we conclude that 0 ≤ tσ/2 Qδm (t) ≤ εσ/2

for 0 ≤ t ≤ ε .

Furthermore, 0 ≤ t Pδm−1 (t)2 =

d  ! 1 − Qδm (t) 1 − Qδm (t) ≤  Qδm (0) t dt /0 1 . /0 . 1 ≤1

d = − dt Qδm (s)

for some s ∈ [0, ε]. Thus we have proven the basic estimate δ √ x∗ − xm,δ X ≤ E εσ/2 + (1 + r) √ + q δ ε

for

0 0, which depends only on σ and r. This case is finished. Case II: ε∗ > 1/q. In this case, we substitute ε = 1/q in (B2) and conclude that √ √ σ/2 x∗ − xm,δ X ≤ E q −σ/2 + (r + 2) q δ ≤ E ε∗ + (r + 2) q δ  σ/(σ+1) r+1 √ ≤ δ σ/(σ+1) E 1/(σ+1) + (r + 2) q δ . σ d δ  Qm(δ) (0). Until now, we have It remains to estimate the quantity q = qm =  dt not used the stopping rule. We will now use this rule to prove the estimate  2/(σ+1) E qm ≤ c (B3) δ for some c > 0, which depends only on σ and r. Analogously to qm , we define  d δ Qm(δ)−1 (0). By the previous lemma, we already know that qm−1 :=  dt   rδ < y δ − Kxm(δ)−1,δ Y ≤ δ + (1 + σ)(σ+1)/2

E (σ+1)/2 qm−1

;

that is,

(1 + σ)(1+σ)/2 E . r−1 δ We have to prove such an estimate for m instead of m − 1. (σ+1)/2

qm−1

≤

(B4)

Choose T > 1 and ρ∗ ∈ (0, 1) with T qm−1 /ρ∗ . From λj,m ≥ λj−1,m−1 for all j = 2, . . . , m, we conclude that   m m−1   d  1 1 1 1 qm =  Qδm (0) = ≤ + = + qm−1 . dt λ λ1,m λ λ1,m j=1 j,m j=1 j,m−1 This implies that qm−1 ≤ ρ∗ qm ≤

ρ∗ + ρ∗ qm−1 ; λ1,m

376

Appendix B

that is, qm−1 ≤

ρ∗ . (1 − ρ∗ ) λ1,m

Finally, we need 1 λ2,m

≤

1 λ1,m−1

≤

m−1 

1

j=1

λj,m−1

= qm−1 ≤

1 ρ∗ . 1 − ρ∗ λ1,m

Now we set ε := T λ1,m . Then ρ∗ λ2,m ≤ 2 λ2,m . 1 − ρ∗ Define the polynomial φ ∈ Pm−1 by  −1  m  9 t t δ φ(t) := Qm (t) 1 − = 1− . λ1,m λj,m j=2 ε ≤ T

For t ≤ ε and j ≥ 2, we note that t ε 1 ≥ 1 − ≥ 1 − ≥ −1 ; λj,m λ2,m that is, |φ(t)| ≤ 1 for all 0 ≤ t ≤ ε . For t ≥ ε, we conclude that     1 − t  = t − λ1,m ≥ ε − 1 = T − 1 ;  λ1,m  λ1,m λ1,m that is,

 1  δ Qm (t) for all t ≥ ε . T −1 Since φ(0) = 1, we can apply Lemma 2.22. Using the projector Mε , we conclude that |φ(t)| ≤

rδ